Diffusers documentation
Image-to-image
Image-to-image
The Stable Diffusion model can also be applied to image-to-image generation by passing a text prompt and an initial image to condition the generation of new images.
The StableDiffusionImg2ImgPipeline uses the diffusion-denoising mechanism proposed in SDEdit: Guided Image Synthesis and Editing with Stochastic Differential Equations by Chenlin Meng, Yutong He, Yang Song, Jiaming Song, Jiajun Wu, Jun-Yan Zhu, Stefano Ermon.
The abstract from the paper is:
Guided image synthesis enables everyday users to create and edit photo-realistic images with minimum effort. The key challenge is balancing faithfulness to the user input (e.g., hand-drawn colored strokes) and realism of the synthesized image. Existing GAN-based methods attempt to achieve such balance using either conditional GANs or GAN inversions, which are challenging and often require additional training data or loss functions for individual applications. To address these issues, we introduce a new image synthesis and editing method, Stochastic Differential Editing (SDEdit), based on a diffusion model generative prior, which synthesizes realistic images by iteratively denoising through a stochastic differential equation (SDE). Given an input image with user guide of any type, SDEdit first adds noise to the input, then subsequently denoises the resulting image through the SDE prior to increase its realism. SDEdit does not require task-specific training or inversions and can naturally achieve the balance between realism and faithfulness. SDEdit significantly outperforms state-of-the-art GAN-based methods by up to 98.09% on realism and 91.72% on overall satisfaction scores, according to a human perception study, on multiple tasks, including stroke-based image synthesis and editing as well as image compositing.
Make sure to check out the Stable Diffusion Tips section to learn how to explore the tradeoff between scheduler speed and quality, and how to reuse pipeline components efficiently!
StableDiffusionImg2ImgPipeline
class diffusers.StableDiffusionImg2ImgPipeline
< source >( vae: AutoencoderKL text_encoder: CLIPTextModel tokenizer: CLIPTokenizer unet: UNet2DConditionModel scheduler: KarrasDiffusionSchedulers safety_checker: StableDiffusionSafetyChecker feature_extractor: CLIPImageProcessor requires_safety_checker: bool = True )
Parameters
- vae (AutoencoderKL) — Variational Auto-Encoder (VAE) model to encode and decode images to and from latent representations.
- text_encoder (CLIPTextModel) — Frozen text-encoder (clip-vit-large-patch14).
-
tokenizer (CLIPTokenizer) —
A
CLIPTokenizerto tokenize text. -
unet (UNet2DConditionModel) —
A
UNet2DConditionModelto denoise the encoded image latents. -
scheduler (SchedulerMixin) —
A scheduler to be used in combination with
unetto denoise the encoded image latents. Can be one of DDIMScheduler, LMSDiscreteScheduler, or PNDMScheduler. -
safety_checker (
StableDiffusionSafetyChecker) — Classification module that estimates whether generated images could be considered offensive or harmful. Please refer to the model card for more details about a model’s potential harms. -
feature_extractor (CLIPImageProcessor) —
A
CLIPImageProcessorto extract features from generated images; used as inputs to thesafety_checker.
Pipeline for text-guided image-to-image generation using Stable Diffusion.
This model inherits from DiffusionPipeline. Check the superclass documentation for the generic methods implemented for all pipelines (downloading, saving, running on a particular device, etc.).
The pipeline also inherits the following loading methods:
- load_textual_inversion() for loading textual inversion embeddings
- load_lora_weights() for loading LoRA weights
- save_lora_weights() for saving LoRA weights
- from_single_file() for loading
.ckptfiles
__call__
< source >(
prompt: typing.Union[str, typing.List[str]] = None
image: typing.Union[torch.FloatTensor, PIL.Image.Image, numpy.ndarray, typing.List[torch.FloatTensor], typing.List[PIL.Image.Image], typing.List[numpy.ndarray]] = None
strength: float = 0.8
num_inference_steps: typing.Optional[int] = 50
guidance_scale: typing.Optional[float] = 7.5
negative_prompt: typing.Union[typing.List[str], str, NoneType] = None
num_images_per_prompt: typing.Optional[int] = 1
eta: typing.Optional[float] = 0.0
generator: typing.Union[torch._C.Generator, typing.List[torch._C.Generator], NoneType] = None
prompt_embeds: typing.Optional[torch.FloatTensor] = None
negative_prompt_embeds: typing.Optional[torch.FloatTensor] = None
output_type: typing.Optional[str] = 'pil'
return_dict: bool = True
callback: typing.Union[typing.Callable[[int, int, torch.FloatTensor], NoneType], NoneType] = None
callback_steps: int = 1
cross_attention_kwargs: typing.Union[typing.Dict[str, typing.Any], NoneType] = None
)
→
StableDiffusionPipelineOutput or tuple
Parameters
-
prompt (
strorList[str], optional) — The prompt or prompts to guide image generation. If not defined, you need to passprompt_embeds. -
image (
torch.FloatTensor,PIL.Image.Image,np.ndarray,List[torch.FloatTensor],List[PIL.Image.Image], orList[np.ndarray]) —Imageor tensor representing an image batch to be used as the starting point. Can also accept image latents asimage, but if passing latents directly it is not encoded again. -
strength (
float, optional, defaults to 0.8) — Indicates extent to transform the referenceimage. Must be between 0 and 1.imageis used as a starting point and more noise is added the higher thestrength. The number of denoising steps depends on the amount of noise initially added. Whenstrengthis 1, added noise is maximum and the denoising process runs for the full number of iterations specified innum_inference_steps. A value of 1 essentially ignoresimage. -
num_inference_steps (
int, optional, defaults to 50) — The number of denoising steps. More denoising steps usually lead to a higher quality image at the expense of slower inference. This parameter is modulated bystrength. -
guidance_scale (
float, optional, defaults to 7.5) — A higher guidance scale value encourages the model to generate images closely linked to the textpromptat the expense of lower image quality. Guidance scale is enabled whenguidance_scale > 1. -
negative_prompt (
strorList[str], optional) — The prompt or prompts to guide what to not include in image generation. If not defined, you need to passnegative_prompt_embedsinstead. Ignored when not using guidance (guidance_scale < 1). -
num_images_per_prompt (
int, optional, defaults to 1) — The number of images to generate per prompt. -
eta (
float, optional, defaults to 0.0) — Corresponds to parameter eta (η) from the DDIM paper. Only applies to the DDIMScheduler, and is ignored in other schedulers. -
generator (
torch.GeneratororList[torch.Generator], optional) — Atorch.Generatorto make generation deterministic. -
prompt_embeds (
torch.FloatTensor, optional) — Pre-generated text embeddings. Can be used to easily tweak text inputs (prompt weighting). If not provided, text embeddings are generated from thepromptinput argument. -
negative_prompt_embeds (
torch.FloatTensor, optional) — Pre-generated negative text embeddings. Can be used to easily tweak text inputs (prompt weighting). If not provided,negative_prompt_embedsare generated from thenegative_promptinput argument. -
output_type (
str, optional, defaults to"pil") — The output format of the generated image. Choose betweenPIL.Imageornp.array. -
return_dict (
bool, optional, defaults toTrue) — Whether or not to return a StableDiffusionPipelineOutput instead of a plain tuple. -
callback (
Callable, optional) — A function that calls everycallback_stepssteps during inference. The function is called with the following arguments:callback(step: int, timestep: int, latents: torch.FloatTensor). -
callback_steps (
int, optional, defaults to 1) — The frequency at which thecallbackfunction is called. If not specified, the callback is called at every step. -
cross_attention_kwargs (
dict, optional) — A kwargs dictionary that if specified is passed along to theAttentionProcessoras defined inself.processor.
Returns
StableDiffusionPipelineOutput or tuple
If return_dict is True, StableDiffusionPipelineOutput is returned,
otherwise a tuple is returned where the first element is a list with the generated images and the
second element is a list of bools indicating whether the corresponding generated image contains
“not-safe-for-work” (nsfw) content.
The call function to the pipeline for generation.
Examples:
>>> import requests
>>> import torch
>>> from PIL import Image
>>> from io import BytesIO
>>> from diffusers import StableDiffusionImg2ImgPipeline
>>> device = "cuda"
>>> model_id_or_path = "runwayml/stable-diffusion-v1-5"
>>> pipe = StableDiffusionImg2ImgPipeline.from_pretrained(model_id_or_path, torch_dtype=torch.float16)
>>> pipe = pipe.to(device)
>>> url = "https://raw.githubusercontent.com/CompVis/stable-diffusion/main/assets/stable-samples/img2img/sketch-mountains-input.jpg"
>>> response = requests.get(url)
>>> init_image = Image.open(BytesIO(response.content)).convert("RGB")
>>> init_image = init_image.resize((768, 512))
>>> prompt = "A fantasy landscape, trending on artstation"
>>> images = pipe(prompt=prompt, image=init_image, strength=0.75, guidance_scale=7.5).images
>>> images[0].save("fantasy_landscape.png")enable_attention_slicing
< source >( slice_size: typing.Union[str, int, NoneType] = 'auto' )
Parameters
-
slice_size (
strorint, optional, defaults to"auto") — When"auto", halves the input to the attention heads, so attention will be computed in two steps. If"max", maximum amount of memory will be saved by running only one slice at a time. If a number is provided, uses as many slices asattention_head_dim // slice_size. In this case,attention_head_dimmust be a multiple ofslice_size.
Enable sliced attention computation. When this option is enabled, the attention module splits the input tensor in slices to compute attention in several steps. For more than one attention head, the computation is performed sequentially over each head. This is useful to save some memory in exchange for a small speed decrease.
⚠️ Don’t enable attention slicing if you’re already using scaled_dot_product_attention (SDPA) from PyTorch
2.0 or xFormers. These attention computations are already very memory efficient so you won’t need to enable
this function. If you enable attention slicing with SDPA or xFormers, it can lead to serious slow downs!
Examples:
>>> import torch
>>> from diffusers import StableDiffusionPipeline
>>> pipe = StableDiffusionPipeline.from_pretrained(
... "runwayml/stable-diffusion-v1-5",
... torch_dtype=torch.float16,
... use_safetensors=True,
... )
>>> prompt = "a photo of an astronaut riding a horse on mars"
>>> pipe.enable_attention_slicing()
>>> image = pipe(prompt).images[0]Disable sliced attention computation. If enable_attention_slicing was previously called, attention is
computed in one step.
enable_xformers_memory_efficient_attention
< source >( attention_op: typing.Optional[typing.Callable] = None )
Parameters
-
attention_op (
Callable, optional) — Override the defaultNoneoperator for use asopargument to thememory_efficient_attention()function of xFormers.
Enable memory efficient attention from xFormers. When this option is enabled, you should observe lower GPU memory usage and a potential speed up during inference. Speed up during training is not guaranteed.
⚠️ When memory efficient attention and sliced attention are both enabled, memory efficient attention takes precedent.
Examples:
>>> import torch
>>> from diffusers import DiffusionPipeline
>>> from xformers.ops import MemoryEfficientAttentionFlashAttentionOp
>>> pipe = DiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-2-1", torch_dtype=torch.float16)
>>> pipe = pipe.to("cuda")
>>> pipe.enable_xformers_memory_efficient_attention(attention_op=MemoryEfficientAttentionFlashAttentionOp)
>>> # Workaround for not accepting attention shape using VAE for Flash Attention
>>> pipe.vae.enable_xformers_memory_efficient_attention(attention_op=None)Disable memory efficient attention from xFormers.
load_textual_inversion
< source >( pretrained_model_name_or_path: typing.Union[str, typing.List[str], typing.Dict[str, torch.Tensor], typing.List[typing.Dict[str, torch.Tensor]]] token: typing.Union[str, typing.List[str], NoneType] = None **kwargs )
Parameters
-
pretrained_model_name_or_path (
stroros.PathLikeorList[str or os.PathLike]orDictorList[Dict]) — Can be either one of the following or a list of them:- A string, the model id (for example
sd-concepts-library/low-poly-hd-logos-icons) of a pretrained model hosted on the Hub. - A path to a directory (for example
./my_text_inversion_directory/) containing the textual inversion weights. - A path to a file (for example
./my_text_inversions.pt) containing textual inversion weights. - A torch state dict.
- A string, the model id (for example
-
token (
strorList[str], optional) — Override the token to use for the textual inversion weights. Ifpretrained_model_name_or_pathis a list, thentokenmust also be a list of equal length. -
weight_name (
str, optional) — Name of a custom weight file. This should be used when:- The saved textual inversion file is in 🤗 Diffusers format, but was saved under a specific weight
name such as
text_inv.bin. - The saved textual inversion file is in the Automatic1111 format.
- The saved textual inversion file is in 🤗 Diffusers format, but was saved under a specific weight
name such as
-
cache_dir (
Union[str, os.PathLike], optional) — Path to a directory where a downloaded pretrained model configuration is cached if the standard cache is not used. -
force_download (
bool, optional, defaults toFalse) — Whether or not to force the (re-)download of the model weights and configuration files, overriding the cached versions if they exist. -
resume_download (
bool, optional, defaults toFalse) — Whether or not to resume downloading the model weights and configuration files. If set toFalse, any incompletely downloaded files are deleted. -
proxies (
Dict[str, str], optional) — A dictionary of proxy servers to use by protocol or endpoint, for example,{'http': 'foo.bar:3128', 'http://hostname': 'foo.bar:4012'}. The proxies are used on each request. -
local_files_only (
bool, optional, defaults toFalse) — Whether to only load local model weights and configuration files or not. If set toTrue, the model won’t be downloaded from the Hub. -
use_auth_token (
stror bool, optional) — The token to use as HTTP bearer authorization for remote files. IfTrue, the token generated fromdiffusers-cli login(stored in~/.huggingface) is used. -
revision (
str, optional, defaults to"main") — The specific model version to use. It can be a branch name, a tag name, a commit id, or any identifier allowed by Git. -
subfolder (
str, optional, defaults to"") — The subfolder location of a model file within a larger model repository on the Hub or locally. -
mirror (
str, optional) — Mirror source to resolve accessibility issues if you’re downloading a model in China. We do not guarantee the timeliness or safety of the source, and you should refer to the mirror site for more information.
Load textual inversion embeddings into the text encoder of StableDiffusionPipeline (both 🤗 Diffusers and Automatic1111 formats are supported).
Example:
To load a textual inversion embedding vector in 🤗 Diffusers format:
from diffusers import StableDiffusionPipeline
import torch
model_id = "runwayml/stable-diffusion-v1-5"
pipe = StableDiffusionPipeline.from_pretrained(model_id, torch_dtype=torch.float16).to("cuda")
pipe.load_textual_inversion("sd-concepts-library/cat-toy")
prompt = "A <cat-toy> backpack"
image = pipe(prompt, num_inference_steps=50).images[0]
image.save("cat-backpack.png")To load a textual inversion embedding vector in Automatic1111 format, make sure to download the vector first (for example from civitAI) and then load the vector
locally:
from diffusers import StableDiffusionPipeline
import torch
model_id = "runwayml/stable-diffusion-v1-5"
pipe = StableDiffusionPipeline.from_pretrained(model_id, torch_dtype=torch.float16).to("cuda")
pipe.load_textual_inversion("./charturnerv2.pt", token="charturnerv2")
prompt = "charturnerv2, multiple views of the same character in the same outfit, a character turnaround of a woman wearing a black jacket and red shirt, best quality, intricate details."
image = pipe(prompt, num_inference_steps=50).images[0]
image.save("character.png")from_single_file
< source >( pretrained_model_link_or_path **kwargs )
Parameters
-
pretrained_model_link_or_path (
stroros.PathLike, optional) — Can be either:- A link to the
.ckptfile (for example"https://huggingface.co/<repo_id>/blob/main/<path_to_file>.ckpt") on the Hub. - A path to a file containing all pipeline weights.
- A link to the
-
torch_dtype (
strortorch.dtype, optional) — Override the defaulttorch.dtypeand load the model with another dtype. If"auto"is passed, the dtype is automatically derived from the model’s weights. -
force_download (
bool, optional, defaults toFalse) — Whether or not to force the (re-)download of the model weights and configuration files, overriding the cached versions if they exist. -
cache_dir (
Union[str, os.PathLike], optional) — Path to a directory where a downloaded pretrained model configuration is cached if the standard cache is not used. -
resume_download (
bool, optional, defaults toFalse) — Whether or not to resume downloading the model weights and configuration files. If set toFalse, any incompletely downloaded files are deleted. -
proxies (
Dict[str, str], optional) — A dictionary of proxy servers to use by protocol or endpoint, for example,{'http': 'foo.bar:3128', 'http://hostname': 'foo.bar:4012'}. The proxies are used on each request. -
local_files_only (
bool, optional, defaults toFalse) — Whether to only load local model weights and configuration files or not. If set toTrue, the model won’t be downloaded from the Hub. -
use_auth_token (
stror bool, optional) — The token to use as HTTP bearer authorization for remote files. IfTrue, the token generated fromdiffusers-cli login(stored in~/.huggingface) is used. -
revision (
str, optional, defaults to"main") — The specific model version to use. It can be a branch name, a tag name, a commit id, or any identifier allowed by Git. -
use_safetensors (
bool, optional, defaults toNone) — If set toNone, the safetensors weights are downloaded if they’re available and if the safetensors library is installed. If set toTrue, the model is forcibly loaded from safetensors weights. If set toFalse, safetensors weights are not loaded. -
extract_ema (
bool, optional, defaults toFalse) — Whether to extract the EMA weights or not. PassTrueto extract the EMA weights which usually yield higher quality images for inference. Non-EMA weights are usually better for continuing finetuning. -
upcast_attention (
bool, optional, defaults toNone) — Whether the attention computation should always be upcasted. -
image_size (
int, optional, defaults to 512) — The image size the model was trained on. Use 512 for all Stable Diffusion v1 models and the Stable Diffusion v2 base model. Use 768 for Stable Diffusion v2. -
prediction_type (
str, optional) — The prediction type the model was trained on. Use'epsilon'for all Stable Diffusion v1 models and the Stable Diffusion v2 base model. Use'v_prediction'for Stable Diffusion v2. -
num_in_channels (
int, optional, defaults toNone) — The number of input channels. IfNone, it is automatically inferred. -
scheduler_type (
str, optional, defaults to"pndm") — Type of scheduler to use. Should be one of["pndm", "lms", "heun", "euler", "euler-ancestral", "dpm", "ddim"]. -
load_safety_checker (
bool, optional, defaults toTrue) — Whether to load the safety checker or not. -
text_encoder (CLIPTextModel, optional, defaults to
None) — An instance ofCLIPTextModelto use, specifically the clip-vit-large-patch14 variant. If this parameter isNone, the function loads a new instance ofCLIPTextModelby itself if needed. -
vae (
AutoencoderKL, optional, defaults toNone) — Variational Auto-Encoder (VAE) Model to encode and decode images to and from latent representations. If this parameter isNone, the function will load a new instance of [CLIP] by itself, if needed. -
tokenizer (CLIPTokenizer, optional, defaults to
None) — An instance ofCLIPTokenizerto use. If this parameter isNone, the function loads a new instance ofCLIPTokenizerby itself if needed. -
kwargs (remaining dictionary of keyword arguments, optional) —
Can be used to overwrite load and saveable variables (for example the pipeline components of the
specific pipeline class). The overwritten components are directly passed to the pipelines
__init__method. See example below for more information.
Instantiate a DiffusionPipeline from pretrained pipeline weights saved in the .ckpt or .safetensors
format. The pipeline is set in evaluation mode (model.eval()) by default.
Examples:
>>> from diffusers import StableDiffusionPipeline
>>> # Download pipeline from huggingface.co and cache.
>>> pipeline = StableDiffusionPipeline.from_single_file(
... "https://huggingface.co/WarriorMama777/OrangeMixs/blob/main/Models/AbyssOrangeMix/AbyssOrangeMix.safetensors"
... )
>>> # Download pipeline from local file
>>> # file is downloaded under ./v1-5-pruned-emaonly.ckpt
>>> pipeline = StableDiffusionPipeline.from_single_file("./v1-5-pruned-emaonly")
>>> # Enable float16 and move to GPU
>>> pipeline = StableDiffusionPipeline.from_single_file(
... "https://huggingface.co/runwayml/stable-diffusion-v1-5/blob/main/v1-5-pruned-emaonly.ckpt",
... torch_dtype=torch.float16,
... )
>>> pipeline.to("cuda")load_lora_weights
< source >( pretrained_model_name_or_path_or_dict: typing.Union[str, typing.Dict[str, torch.Tensor]] **kwargs )
Parameters
-
pretrained_model_name_or_path_or_dict (
stroros.PathLikeordict) — See lora_state_dict(). -
kwargs (
dict, optional) — See lora_state_dict().
Load LoRA weights specified in pretrained_model_name_or_path_or_dict into self.unet and
self.text_encoder.
All kwargs are forwarded to self.lora_state_dict.
See lora_state_dict() for more details on how the state dict is loaded.
See load_lora_into_unet() for more details on how the state dict is loaded into
self.unet.
See load_lora_into_text_encoder() for more details on how the state dict is loaded
into self.text_encoder.
save_lora_weights
< source >( save_directory: typing.Union[str, os.PathLike] unet_lora_layers: typing.Dict[str, typing.Union[torch.nn.modules.module.Module, torch.Tensor]] = None text_encoder_lora_layers: typing.Dict[str, torch.nn.modules.module.Module] = None is_main_process: bool = True weight_name: str = None save_function: typing.Callable = None safe_serialization: bool = True )
Parameters
-
save_directory (
stroros.PathLike) — Directory to save LoRA parameters to. Will be created if it doesn’t exist. -
unet_lora_layers (
Dict[str, torch.nn.Module]orDict[str, torch.Tensor]) — State dict of the LoRA layers corresponding to theunet. -
text_encoder_lora_layers (
Dict[str, torch.nn.Module]orDict[str, torch.Tensor]) — State dict of the LoRA layers corresponding to thetext_encoder. Must explicitly pass the text encoder LoRA state dict because it comes from 🤗 Transformers. -
is_main_process (
bool, optional, defaults toTrue) — Whether the process calling this is the main process or not. Useful during distributed training and you need to call this function on all processes. In this case, setis_main_process=Trueonly on the main process to avoid race conditions. -
save_function (
Callable) — The function to use to save the state dictionary. Useful during distributed training when you need to replacetorch.savewith another method. Can be configured with the environment variableDIFFUSERS_SAVE_MODE. -
safe_serialization (
bool, optional, defaults toTrue) — Whether to save the model usingsafetensorsor the traditional PyTorch way withpickle.
Save the LoRA parameters corresponding to the UNet and text encoder.
Offload all models to CPU to reduce memory usage with a low impact on performance. Moves one whole model at a
time to the GPU when its forward method is called, and the model remains in GPU until the next model runs.
Memory savings are lower than using enable_sequential_cpu_offload, but performance is much better due to the
iterative execution of the unet.
StableDiffusionPipelineOutput
class diffusers.pipelines.stable_diffusion.StableDiffusionPipelineOutput
< source >( images: typing.Union[typing.List[PIL.Image.Image], numpy.ndarray] nsfw_content_detected: typing.Optional[typing.List[bool]] )
Parameters
-
images (
List[PIL.Image.Image]ornp.ndarray) — List of denoised PIL images of lengthbatch_sizeor NumPy array of shape(batch_size, height, width, num_channels). -
nsfw_content_detected (
List[bool]) — List indicating whether the corresponding generated image contains “not-safe-for-work” (nsfw) content orNoneif safety checking could not be performed.
Output class for Stable Diffusion pipelines.
FlaxStableDiffusionImg2ImgPipeline
class diffusers.FlaxStableDiffusionImg2ImgPipeline
< source >( vae: FlaxAutoencoderKL text_encoder: FlaxCLIPTextModel tokenizer: CLIPTokenizer unet: FlaxUNet2DConditionModel scheduler: typing.Union[diffusers.schedulers.scheduling_ddim_flax.FlaxDDIMScheduler, diffusers.schedulers.scheduling_pndm_flax.FlaxPNDMScheduler, diffusers.schedulers.scheduling_lms_discrete_flax.FlaxLMSDiscreteScheduler, diffusers.schedulers.scheduling_dpmsolver_multistep_flax.FlaxDPMSolverMultistepScheduler] safety_checker: FlaxStableDiffusionSafetyChecker feature_extractor: CLIPImageProcessor dtype: dtype = <class 'jax.numpy.float32'> )
Parameters
- vae (FlaxAutoencoderKL) — Variational Auto-Encoder (VAE) model to encode and decode images to and from latent representations.
- text_encoder (FlaxCLIPTextModel) — Frozen text-encoder (clip-vit-large-patch14).
-
tokenizer (CLIPTokenizer) —
A
CLIPTokenizerto tokenize text. -
unet (FlaxUNet2DConditionModel) —
A
FlaxUNet2DConditionModelto denoise the encoded image latents. -
scheduler (SchedulerMixin) —
A scheduler to be used in combination with
unetto denoise the encoded image latents. Can be one ofFlaxDDIMScheduler,FlaxLMSDiscreteScheduler,FlaxPNDMScheduler, orFlaxDPMSolverMultistepScheduler. -
safety_checker (
FlaxStableDiffusionSafetyChecker) — Classification module that estimates whether generated images could be considered offensive or harmful. Please refer to the model card for more details about a model’s potential harms. -
feature_extractor (CLIPImageProcessor) —
A
CLIPImageProcessorto extract features from generated images; used as inputs to thesafety_checker.
Flax-based pipeline for text-guided image-to-image generation using Stable Diffusion.
This model inherits from FlaxDiffusionPipeline. Check the superclass documentation for the generic methods implemented for all pipelines (downloading, saving, running on a particular device, etc.).
__call__
< source >(
prompt_ids: array
image: array
params: typing.Union[typing.Dict, flax.core.frozen_dict.FrozenDict]
prng_seed: PRNGKeyArray
strength: float = 0.8
num_inference_steps: int = 50
height: typing.Optional[int] = None
width: typing.Optional[int] = None
guidance_scale: typing.Union[float, array] = 7.5
noise: array = None
neg_prompt_ids: array = None
return_dict: bool = True
jit: bool = False
)
→
FlaxStableDiffusionPipelineOutput or tuple
Parameters
-
prompt_ids (
jnp.array) — The prompt or prompts to guide image generation. -
image (
jnp.array) — Array representing an image batch to be used as the starting point. -
params (
DictorFrozenDict) — Dictionary containing the model parameters/weights. -
prng_seed (
jax.random.KeyArrayorjax.Array) — Array containing random number generator key. -
strength (
float, optional, defaults to 0.8) — Indicates extent to transform the referenceimage. Must be between 0 and 1.imageis used as a starting point and more noise is added the higher thestrength. The number of denoising steps depends on the amount of noise initially added. Whenstrengthis 1, added noise is maximum and the denoising process runs for the full number of iterations specified innum_inference_steps. A value of 1 essentially ignoresimage. -
num_inference_steps (
int, optional, defaults to 50) — The number of denoising steps. More denoising steps usually lead to a higher quality image at the expense of slower inference. This parameter is modulated bystrength. -
height (
int, optional, defaults toself.unet.config.sample_size * self.vae_scale_factor) — The height in pixels of the generated image. -
width (
int, optional, defaults toself.unet.config.sample_size * self.vae_scale_factor) — The width in pixels of the generated image. -
guidance_scale (
float, optional, defaults to 7.5) — A higher guidance scale value encourages the model to generate images closely linked to the textpromptat the expense of lower image quality. Guidance scale is enabled whenguidance_scale > 1. -
noise (
jnp.array, optional) — Pre-generated noisy latents sampled from a Gaussian distribution to be used as inputs for image generation. Can be used to tweak the same generation with different prompts. The array is generated by sampling using the supplied randomgenerator. -
return_dict (
bool, optional, defaults toTrue) — Whether or not to return a FlaxStableDiffusionPipelineOutput instead of a plain tuple. -
jit (
bool, defaults toFalse) — Whether to runpmapversions of the generation and safety scoring functions.This argument exists because
__call__is not yet end-to-end pmap-able. It will be removed in a future release.
Returns
FlaxStableDiffusionPipelineOutput or tuple
If return_dict is True, FlaxStableDiffusionPipelineOutput is
returned, otherwise a tuple is returned where the first element is a list with the generated images
and the second element is a list of bools indicating whether the corresponding generated image
contains “not-safe-for-work” (nsfw) content.
The call function to the pipeline for generation.
Examples:
>>> import jax
>>> import numpy as np
>>> import jax.numpy as jnp
>>> from flax.jax_utils import replicate
>>> from flax.training.common_utils import shard
>>> import requests
>>> from io import BytesIO
>>> from PIL import Image
>>> from diffusers import FlaxStableDiffusionImg2ImgPipeline
>>> def create_key(seed=0):
... return jax.random.PRNGKey(seed)
>>> rng = create_key(0)
>>> url = "https://raw.githubusercontent.com/CompVis/stable-diffusion/main/assets/stable-samples/img2img/sketch-mountains-input.jpg"
>>> response = requests.get(url)
>>> init_img = Image.open(BytesIO(response.content)).convert("RGB")
>>> init_img = init_img.resize((768, 512))
>>> prompts = "A fantasy landscape, trending on artstation"
>>> pipeline, params = FlaxStableDiffusionImg2ImgPipeline.from_pretrained(
... "CompVis/stable-diffusion-v1-4",
... revision="flax",
... dtype=jnp.bfloat16,
... )
>>> num_samples = jax.device_count()
>>> rng = jax.random.split(rng, jax.device_count())
>>> prompt_ids, processed_image = pipeline.prepare_inputs(
... prompt=[prompts] * num_samples, image=[init_img] * num_samples
... )
>>> p_params = replicate(params)
>>> prompt_ids = shard(prompt_ids)
>>> processed_image = shard(processed_image)
>>> output = pipeline(
... prompt_ids=prompt_ids,
... image=processed_image,
... params=p_params,
... prng_seed=rng,
... strength=0.75,
... num_inference_steps=50,
... jit=True,
... height=512,
... width=768,
... ).images
>>> output_images = pipeline.numpy_to_pil(np.asarray(output.reshape((num_samples,) + output.shape[-3:])))FlaxStableDiffusionPipelineOutput
class diffusers.pipelines.stable_diffusion.FlaxStableDiffusionPipelineOutput
< source >( images: ndarray nsfw_content_detected: typing.List[bool] )
Output class for Flax-based Stable Diffusion pipelines.
“Returns a new object replacing the specified fields with new values.