text
stringlengths
2
806k
meta
dict
1. Field of the Invention The present invention relates to an evoked response detector for a heart stimulator for determining evoked response in the presence of polarization, and to a heart stimulator having such an evoked response detector incorporated therein. 2. Description of the Prior Art Cardiac stimulators are known which have a pulse generator devised for producing stimulation pulses of varying amplitudes, and a lead adopted to be introduced into the heart of a patient and connected to the pulse generator for delivering stimulation pulses to the heart, and an evoked response detector having measuring and memory means for measuring and storing the electrode signal picked up by the lead in response to delivered stimulation pulses, wherein at least one of said stimulation pulses has a sufficiently high amplitude for obtaining capture. To reduce the energy consumption of heart stimulators a so called AUTOCAPTURE.TM. function is used to maintain the energy of the stimulation pulses at a level just above that which is needed to effectuate capture, cf. e.g. U.S. Pat. No. 5,458,623. A reliable detection of the evoked response, which then is necessary, is, however, not a simple matter, especially when it is desired to sense the evoked response with the same electrode as the one delivering the stimulation pulse. This is because of the fact that the evoked response potential could be small in amplitude compared to the residual polarization charge. The residual charge decays exponentially but tends to dominate the evoked potential for several hundreds of milliseconds after the stimulation. If the polarization is too high, it could be wrongly interpreted by the evoked response detector as a capture, i.e. contraction of the heart. The AUTOCAPTURE.TM. algorithm could then by mistake adjust the output amplitude of the stimulation pulse to a value below the actual capture level, which will result in no capture. If the electrode surface of the electrode lead in use has significant polarization this could consequently disturb the AUTOCAPTURE.TM. function and result in loss of capture. To guarantee a safe and reliable detection of evoked response it is thus desirable to use leads having electrode surfaces with low polarization. Several attempts have been made to solve the lead polarization problems in connection with evoked response detection. Thus U.S. Pat. No. 5,417,718 discloses a system for maintaining capture wherein electrical post-stimulus signal of the heart, following delivery of a stimulation pulse, is compared to a polarization template, determined during a capture verification test. A prescribed difference between the polarization template and the post-stimulus signal indicates capture. Otherwise loss of capture is presumed and the stimulation energy is increased a predetermined amount to obtain capture. There is mostly at least one significant slope in the bipolar measured IEGM signal, which makes it possible to discriminate the evoked response signal from slowly varying signals such as polarization signals. Thus in U.S. Pat. No. 5,431,693 a method of verifying capture of the heart by a cardiac pacemaker is described by observing that the non-capture potential is exponential in form and the evoked capture potential, while generally exponential in form, has one or more small-amplitude perturbations superimposed on the exponential wave form. These perturbations are enhanced for ease of detection by processing the wave form signal by differentiation to form the second derivative of the evoked response signal for analysis for the evoked response detection. Unipolar detection of evoked response signals is, however, not possible using this technique. Abrupt slope changes or superimposed small-amplitude perturbations are leveled out if the measurements are made over a longer distance from the electrode to the stimulator casing. Experiments have now shown that the evoked response signal amplitude is fairly constant, independent of the stimulation pulse amplitude, i.e. the evoked response signal amplitude does not vary with the amplitude of the stimulation pulse (provided that the stimulation amplitude is above the capture threshold) Further, it has been found that the electrode polarization is approximately linearly dependent on the stimulation pulse amplitude for a constant pulse duration. Experimental results are presented in greater detail below in connection with the description of FIGS. 1-3.
{ "pile_set_name": "USPTO Backgrounds" }
1. Field of Invention The present invention relates, generally, to automotive powertrain systems and, more specifically, to a wheel-end disconnect assembly for powertrain systems. 2. Description of the Related Art Conventional automotive vehicles known in the art include a powertrain system in rotational communication with one or more drivelines. Typically, the vehicle includes a pair of drivelines, each defined by a respective pair of opposing wheels. The powertrain system includes a propulsion system adapted to generate and selectively translate rotational torque to one or more of the wheels so as to drive the vehicle. To that end, in conventional automotive powertrain systems, the propulsion system is typically realized as an internal combustion engine in rotational communication with a transmission. The engine generates rotational torque which is selectively translated to the transmission which, in turn, translates rotational torque to one or more of the drivelines. The transmission multiplies the rotational speed and torque generated by the engine through a series of predetermined gear sets, whereby changing between gear sets enables the vehicle to travel at different vehicle speeds for a given engine speed. In so-called “four-wheel-drive” or “all-wheel-drive” powertrain systems, both drivelines are used to drive the vehicle. To that end, all wheel drive powertrain systems typically include a transfer case disposed in rotational communication with the transmission and adapted to split rotational torque between the drivelines. The transfer case may be spaced from the transmission, or may be integrated with the transmission. Where the transfer case is spaced from the transmission, a driveshaft is used to translate rotational torque from the transmission to the transfer case. Driveshafts are also typically used to connect the transfer case to each respective driveline. Conventional drivelines are commonly realized by a differential assembly adapted to receive rotational torque from the transfer case and subsequently split rotational torque between opposing wheels. To that end, each driveline also typically includes a pair of continuously-variable joints disposed in torque translating relationship with the differential and each respective opposing wheel. Depending on the specific configuration of the powertrain system, the percentage of torque split between the drivelines may vary. Moreover, depending on the vehicle application, the transfer case and/or driveline(s) may be configured to interrupt rotational torque to one of the drivelines under certain operating conditions. Specifically, the powertrain system may be configured such that the vehicle can be selectively operated in “two-wheel-drive” or in “four-wheel-drive”. Moreover, the powertrain system may be configured to automatically and continuously control how much rotational torque is sent to each driveline. Thus, the powertrain system may be configured to send a higher percentage of available torque to one of the drivelines under certain vehicle operating conditions, and a lower percentage of available torque to the same driveline under different vehicle operating conditions. By way of non-limiting example, the powertrain system may be configured such that 80% of torque is sent to a front driveline and 20% of torque is sent to a rear driveline until there is a loss of traction or wheel spin, whereby the powertrain subsequently adjusts torque split such that 50% of torque is sent to each driveline. Depending on the vehicle application, rotational torque may only be required at both drivelines relatively infrequently. Thus, the vehicle may be designed to operate primarily in “two-wheel-drive” so as to minimize parasitic loss and optimize powertrain system efficiency. Moreover, optimizing how torque is split between drivelines can lead to significant improvements in vehicle efficiency. Thus, in order to decrease parasitic losses in the powertrain system, it is advantageous to selectively disconnect one or more driveshafts and/or continuously-variable joints from rotational communication with the transfer case, transmission, and/or differentials. To that end, rotational disconnects are used to selectively interrupt rotation between powertrain system components, whereby a controller and an actuator are typically used to selectively control the rotational disconnect. The controller energizes the actuator which, in turn, engages the rotational disconnect so as to couple (or, de-couple) the powertrain system components. Each of the components and systems of the type described above must cooperate to effectively and selectively translate rotational torque to the driven wheels of the vehicle. In addition, each of the components and systems must be designed not only to facilitate improved performance and efficiency, but also so as to reduce the cost and complexity of manufacturing vehicles. While powertrain rotational disconnect systems known in the related art have generally performed well for their intended use, there remains a need in the art for a wheel-end rotational disconnect assembly that has superior operational characteristics and a reduced overall packaging size, and, at the same time, that reduces the cost and complexity of manufacturing vehicles that operate with high efficiency under a number of different driving conditions.
{ "pile_set_name": "USPTO Backgrounds" }
1. Statement of the Technical Field The inventive arrangements relate to portable radios, and more particularly to maximizing transmitter power output of portable handheld radio equipment. 2. Description of the Related Art Improvements in technology have led to the development of a new generation of multi-band portable radios. These portable radios include multi-mission tactical radios that are designed to be compatible with the Joint Tactical Radio System (JTRS). The JTRS system is built upon the well known Software Communications Architecture (SCA) which establishes a system by which programmable radios can be quickly modified with software to implement any one of a wide variety of communication link protocols on a variety of frequency bands. These highly flexible radios are capable of operating over a very wide range of frequencies, for communicating voice and data using any one of a variety of modulation schemes. Notwithstanding the highly flexible software defined nature of these radios, there remain certain aspects of the equipment that must be implemented in hardware. This hardware typically includes the transmitter RF power amplifier circuitry. The hardware nature of this circuitry has typically resulted in certain compromises. For example, since the nature of the transmitted signals is programmable, the transmitter power amplifier circuitry must be designed to accommodate any communications protocol, regardless of frequency, waveform, modulation, data rate, and duty cycle. While this approach ensures flexibility, it is necessarily a compromise design. For example, the transmitter output power can be limited to protect the final RF power amplifier under worst case conditions. This is typically implemented at design time by fixing various circuit parameters, such as the voltage applied to the power amplifier, transistor bias current settings, and maximum power output. While necessary to protect the RF power amplifier, these fixed design parameters can limit the maximum transmitter performance achievable under certain circumstances. For example, an RF power amplifier can be designed for a maximum 5 watt output power. This design can protect the power amplifier from damage caused by overheating. Under certain limited circumstances it may be possible to safely increase power amplifier RF power output levels. However, because the RF power amplifier design must be compatible with any communication protocol, the more conservative circuit design must be implemented.
{ "pile_set_name": "USPTO Backgrounds" }
Signal transducer and activator of transcription (STAT) proteins are transcription factors which transduce signals from various extracellular cytokines and growth factors to a nucleus. Seven (7) subtypes of STAT proteins (i.e., STAT1, STAT2, STAT3, STAT4, STAT5a, STAT5b and STAT6) are currently known, and generally they consist of about 750-850 amino acids. In addition, each subtype of STAT proteins contains several conserved domains which play an important role in exhibiting the function of STAT proteins. Specifically, five (5) domains from N-terminus to C-terminus of STAT proteins have been reported including coiled-coiled domain, DNA binding domain, linker domain, SH2 domain and transactivation domain (TAD)). Further, X-ray crystalline structures of STAT1, STAT3, STAT4 and STAT5 have been reported since 1998 (Becker S et al., Nature, 1998, 394; Vinkemeier U et al., Science, 1998, 279; Chen X et al., Cell, 1998, 93; D. Neculai et al., J. Biol. Chem., 2005, 280). In general, receptors to which cytokines and growth factors bind are categorized into Class I and Class II. IL-2, IL-3, IL-5, IL-6, IL-12, G-CSF, GM-CSF, LIF, thrombopoietin, etc., bind to Class I receptors, while INF-α, INF-γ, IL-10, etc., bind to Class II receptors (Schindler C et al., Annu. Rev. Biochem., 1995, 64; Novick D et al., Cell, 1994, 77; Ho AS et al., Proc. Natl. Acad. Sci., 1993, 90). Among them, the cytokine receptors involved in the activation of STAT proteins can be classified depending on their structural forms of extracellular domains into a gp-130 family, an IL-2 family, a growth factor family, an interferon family and a receptor tyrosine kinase family. Interleukin-6 family cytokines are representative multifunctional cytokines which mediate various physiological activities. When interleukin-6 cytokine binds to IL-6 receptor, it attracts gp-130 receptor to form an IL-6-gp-130 receptor complex. At the same time, JAK kinases (JAK1, JAK2, JAK3 and Tyk2) in the cytoplasm are recruited to a cytoplasmic region of gp130 to be phosphorylated and activated. Subsequently, latent cytoplasmic STAT proteins are attracted to a receptor, phosphorylated by JAK kinases and activated. Tyrosine-705 near the SH2 domain located in the C-terminus of STAT proteins is phosphorylated, and the activated tyrosine-705 of each STAT protein monomer binds to the SH2 domain of another monomer in a reciprocal manner, thereby forming a homo- or heterodimer. The dimers are translocalized into a nucleus and bind to a specific DNA binding promoter to promote the transcription. Through its transcription process, various proteins (Myc, Cyclin D1/D2, Bcl-xL, Mcl, survivin, VEGF, HIF-1, immune suppressors, etc.) associated with cell proliferation, survival, angiogenesis and immune evasion are produced (Stark et al., Annu. Rev. Biochem., 1997, 67; Levy et al., Nat. Rev. Mol. Cell Biol., 2002, 3). In particular, STAT3 protein is known to play a crucial role in the acute inflammatory response and the signal transduction pathway of IL-6 and EGF (Akira et al., Cell, 1994, 76; Zhong et al., Science, 1994, 264). According to the recent clinical report, STAT3 protein is constantly activated in patients with solid cancers occurring in prostate, stomach, breast, lung, pancreas, kidney, uterine, ovary, head and neck, etc., and also in patients with blood cancer such as acute and chronic leukemia, multiple myeloma, etc. Further, it has been reported that the survival rate of a patient group with activated STAT3 is remarkably lower than that of a patient group with inactivated STAT3 (Masuda et al., Cancer Res., 2002, 62; Benekli et al., Blood, 2002, 99; Yuichi et al., Int. J. Oncology, 2007, 30). Meanwhile, STAT3 was identified to be an essential factor for the growth and maintenance of murine embryonic stem cells in a study employing a STAT3 knockout mouse model. Also, a study with a tissue-specific STAT3-deficient mouse model, reveals that STAT3 plays an important role in cell growth, apoptosis, and cell motility in a tissue-specific manner (Akira et al., Oncogene 2000, 19). Moreover, since apoptosis by anti-sensing STAT3 was observed in various cancer cell lines, STAT3 is considered as a promising new anticancer target. STAT3 is also considered as a potential target in the treatment of patients with diabetes, immune-related diseases, hepatitis C, macular degeneration, human papillomavirus infection, non-Hodgkin's lymphoma, tuberculosis, etc. In contrast, although STAT1 has identical intracellular response pathways of cytokines and growth factors to those of STAT3, STAT1 increases inflammation and congenital and acquired immunities to inhibit the proliferation of cancer cells or cause pro-apoptotic responses, unlike STAT3 (Valeria Poli et al., Review, Landes Bioscience, 2009). In order to develop STAT3 inhibitors, the following methods can be considered: i) inhibition of the phosphorylation of STAT3 protein by IL-6/gp-130/JAK kinase, ii) inhibition of the dimerization of activated STAT3 protein, and iii) inhibition of the binding of STAT3 dimer to nuclear DNA. Small molecular STAT3 inhibitors are currently under development. Specifically, OPB-31121 and OPB-51602 are under clinical studies on patients with solid cancers or blood cancers by Otsuka Pharmaceutical Co., Ltd. Further, S3I-201 (Siddiquee et al., Proc. Natl. Acad. Sci., 2007, 104), S3I-M2001 (Siddiquee et al., Chem. Biol., 2007, 2), LLL-12 (Lin et al., Neoplasia, 2010, 12), Stattic (Schust et al., Chem. Biol. 2006, 13), STA-21 (Song et al., Proc. Natl. Acad. Sci., 2005, 102), SF-1-066 (Zhang et al., Biochem. Pharm., 2010, 79) and STX-0119 (Matsuno et al., ACS Med. Chem. Lett., 2010, 1), etc. have been tested in a cancer cell growth inhibition experiment and in animal model (in vivo Xenograft model). Furthermore, although peptide compounds mimicking the adjacent amino acid sequence of pY-705 (STAT3) which binds to SH2 domain or the amino acid sequence of gp-130 receptor that JAK kinases bind were studied (Coleman et al., J. Med. Chem., 2005, 48), the development of the peptide compounds has not been successful due to the problems such as solubility and membrane permeability.
{ "pile_set_name": "USPTO Backgrounds" }
A battery is required to power an electronic device when there is no constant source of electrical power or when it is intolerable for the source of electrical power to be interrupted. Whenever a battery-operated device is used, it is desirable to maximize its operating life. One way of doing this is to connect a number of batteries in parallel. One problem with connecting batteries directly in parallel is that current will flow from the battery with the highest voltage to the battery with the lowest voltage. Such a current may cause one or more of the batteries to burst into flame, explode, or both. A diode may be used to prevent current flow from one battery to another. Unfortunately, a diode will also cause an undesirable voltage drop. Such a voltage drop may require the use of a higher voltage battery. To increase the voltage of a battery typically requires an increase in the number of voltaic cells that make up the battery. Therefore, increasing battery voltage typically results in an increase in the size and weight of the battery, both of which may be at a premium in a battery-operated device. Another problem with batteries is determining how many hours of operation are left on a battery. Battery voltage alone is not a good indicator of how much operating time is left on the battery. Typically, the current drain per hour on the battery is measured, the total current capacity of the battery is estimated, and the second number is divided by the first number to arrive at the number of hours of operation for the battery. Such a calculation is difficult and inaccurate if the current drain on the battery is not constant. U.S. Pat. Nos. 3,666,961, entitled "ELECTRICAL POWER SUPPLY"; 3,666,962, entitled "ELECTRICAL POWER SUPPLY"; 5,764,032, entitled "MULTIPLE BATTERY SWITCHOVER CIRCUITS"; and 5,898,291, entitled "BATTERY CELL BYPASS TOPOLOGY," each disclose a device for switching between batteries. However, U.S. Pat. Nos. 3,666,961 and 3,666,962 use high current devices such a solenoid and a current meters which could not maximize the operating life of a battery-operated device as does the present invention. U.S. Pat. Nos. 5,764,032 and 5,898,291 do not provide the functionality as does the present invention. U.S. Pat. Nos. 3,666,961; 3,666,962; 5,764,032, and 5,898,291 are hereby incorporated by reference into the specification of the present invention. U.S. Pat. Nos. 4,622,508, entitled "LITHIUM BATTERY PROTECTION CIRCUIT"; 5,610,495, entitled "CIRCUIT AND METHOD OF MONITORING BATTERY CELLS"; and 5,894,212, entitled "DISCHARGE MONITORING AND ISOLATING SYSTEM FOR BATTERIES," each disclose a device for testing and switching out series-connected individual voltaic cells that make up one battery and does function on multiple batteries connected in parallel as does the present invention. U.S. Pat. Nos. 4,622,508 and 5,610,495 are hereby incorporated by reference into the specification of the present invention.
{ "pile_set_name": "USPTO Backgrounds" }
1. Field of the Invention The present invention is directed to a notepad or notebook computer interface that facilitates input via a pen and, more particularly, to a system that provides interface elements particularly suitable for pen-based input to, and interaction with, electronic documents. 2. Description of the Related Art Making electronic information as handy as paper requires an interaction model which is somewhat different than conventional systems. Documents, the primary form in which people obtain and produce information, are most familiar and manipulatable in the form of paper pages. What is needed is a notepad or notebook computer based system that allows the easy manipulation of document pages. Marking on documents directly, with a pen, provides the most leverage of human motor skills for adding to and interacting with documents. What is needed is a pen based method of marking on and interacting with documents within a notepad or notebook computer. Organizing information for recall is best done according to one's own experience and mental associations, i.e. “when, where I saw it last”, and “what I think it's related to.” What is needed is a system that will organize documents according to the user's experience and mental model. Such a device and its software needs to: render in pages, favoring a portrait aspect ratio and dedicating as much of the screen as possible to displaying the page; support tasks via transitional user interface (U/I) elements which may appear inside or outside or overlapping with the page, and which are controlled directly with the tip of the actual pen (there is no other representation of the pen's location) taking into account the users left or right-handedness, with appearance and behaviors that are both obvious and unobtrusive, easy-to learn and efficient; and support the layering of ink and transitional U/I elements over the content in ways that minimize obstruction of the content and include presentation of such transitional elements in a way which exploits the viewer's ability to resolve overlapping and translucent images. Additionally, the system must support text entry via the pen, but not let the conversion of pen input to text impede the perceived natural flow of events. A notebook computer should behave and allow interaction with it in a manner similar to interaction with a paper notebook. Likewise, a notepad computer should behave and allow interaction with it in a manner similar to interaction with a paper notepad. For example, when flipping pages of paper the user can separate a corner of the page from a stack to make it easier to flip. The position of the user within a stack of pages of a paper is readily visible via looking at the edge of the stack. A paper user can change pens (inks) with the simple motion of picking up another pen. What is needed is an interface that provides such functions simply and intuitively.
{ "pile_set_name": "USPTO Backgrounds" }
1. Field of the Invention The invention relates to aircraft in general and to lifting body aircraft in particular. 2. Prior Art Lifting body aircraft are known in the art. The theory was pioneered by Dr. Alfired J. Eggers, Jr. who discovered that lift could be generated by modifying the shape of a blunt nose cone reentry vehicle. Eggers'theories led to several NASA and U.S. Air Force experimental aircraft in the 1960's and 1970's, namely the M2-F1, the M2-F2, the M2-F3, HL10, X-24A and X-24B. Eggers'lift concept can also be seen in U.S. Pat. No. 3,276,722. All of these craft generate lift with the fuselage of the aircraft only; i.e. without a wing. The Eggers craft can generally be described as conical sections. They have a blunt nose, swept back sides, round bottoms, and generally flat tops. As the crafts move through air, the high profile of their curved lower surfaces causes a significant amount of air to be displaced up and around the body of the vessel. This does two things. First, pressure is being exerted on the air mass at the interface with the hull. In response to this pressure, the air is compressed and it is displaced, up and around the lower surface of the craft. Air has viscosity, so it resists both the compression and the displacement. The greater the speed at which the craft is moving, the greater the resistance of the air. This resistance of the air is transmitted as force to its surroundings. Thus, the air pressure below the craft will be increased. Second, a lift body is designed to move forward with its nose slightly elevated. As the craft moves forward, the craft will displace air, rarefying (reducing the density of) the air above the craft. This will result in a decrease in pressure above the upper surface of the lifting body. The difference between the increased pressure generated on the lower surface and the decreased pressure at the upper surface will result in an overall upward force on the lifting body. This upward force is lift. Although lift bodies were previously known, limited maneuverability and stability of many of the craft in atmospheric flight created significant challenges. Moreover, the inventors are not aware of any that were ever able to takeoff under their own power. Rather, previous lifting bodies were all believed to carried aloft by a winged aircraft and released to perform rocket powered maneuvers and then to return to earth at a glide. Accordingly, a lifting body craft meeting the following objectives is desired.
{ "pile_set_name": "USPTO Backgrounds" }
In total hip replacement surgery, a patient's natural hip is replaced by an acetabular cup component that replaces the acetabular socket and a femoral component that replaces the femoral head. The femoral component of the hip prosthesis includes a generally spherical head, connected via a neck portion, to an elongate stem. The patient's femur is prepared to receive the stem. The proximal end of the femur is resected to expose the medullary canal. This involves resection of at least part of the greater femoral trochanter, and the creation of a cavity that matches the shape of the implant stem. Surgeons may use several different instruments such as an osteotome, rasp, canal probe, and starter broach to initiate the canal. Changing between instruments takes time. Additionally, the number of instruments required for surgery indirectly increases cost of a procedure. There is therefore a need for a multifunctional instrument that combines the functionality of existing instruments. To help ensure proper final orientation of the stem, lateral bias during implant preparation may be preferred. Retraction of the gluteus medius and removal of the lateral cortical bone at the trochanteric fossa helps the surgeon to obtain optimal proximal fit of the stem. This also reduces the risk of undersizing and/or varus placement of the stem. Surgeons currently face several challenges when using traditional osteotomes. For example, the surgeon conventionally positions the osteotome relative to external anatomical landmarks, such as the trochanteric fossa, in an effort to position the osteotome to cut the desired bone leading to the femoral canal. Since every bone is shaped differently, the surgeon may not correctly predict the starting position which would lead to the desired bone removal. This may be further exacerbated by a poor grip between the osteotome and the bone surface, resulting in slippage of the osteotome after initial placement caused by impaction strikes. The surgeon also conventionally angles the osteotome relative to the external anatomical landmarks, such as the femoral leg axis, in an effort to direct the angle that the osteotome resects the bone. Since it is often difficult to visualize the external anatomical landmarks and to hold the osteotome at the desired angle while impacting, the surgeon may not achieve the desired angle of bone removal. The necessary amount of lateral bone may not be fully removed by the first cut. This is particularly true if the surgeon adopts a conservative and iterative bone cutting approach. Thus, further cuts with the osteotome, or removal of the bone with other instruments (e.g., rasp, rotational initiator or broach) may be required. There is therefore a need for an improved osteotome that does not rely on external visual landmarks for positioning. There is also a need for an osteotome having an improved grip on the proximal surface of the femur.
{ "pile_set_name": "USPTO Backgrounds" }
The amount of data that Wireless Wide Area Networks (WWAN) are being required to handle continues to increase. Some reasons for these increases include the accelerated adoption of smart phones and tablets and the emergence of cloud computing applications. Unfortunately, the deployment of additional cellular radio towers (tower) to meet these demands is costly. Additionally, implementation of advanced technologies within wireless cellular networks is also costly, is done over extended periods of time, and involves other inherent constraints. Wireless Local Areas Networks (WLAN) provide an alternative mode of wireless communication. Non-limiting examples of such networks are provided in the specifications of Institute of Electrical and Electronics Engineers (IEEE) 802.11 and IEEE 802.15. However, the localized nature of WLANs prevent them from supporting true mobility in wireless communications. Indeed, the differences between WWANs and WLANs have presented obstacles to the use of one type of network to support the other. Reference will now be made to the exemplary embodiments illustrated, and specific language will be used herein to describe the same. It will nevertheless be understood that no limitation of the scope of the invention is thereby intended.
{ "pile_set_name": "USPTO Backgrounds" }
This invention pertains in general to analog-to-digital converters and in particular to analog-to-digital converters having a very high operating clock frequency, small die size, and low power consumption and methods of stabilizing the same against drift. Conventional high-speed analog-to-digital converters (xe2x80x9cADCsxe2x80x9d) commonly employ a full flash architecture in which the analog-to-digital conversion is done in parallel by using approximately 2n voltage comparators. FIG. 1 illustrates a conventional full flash ADC 100 including an input voltage 110, a reference voltage 112, a number of resistors, of which resistor 114 is representative, a number of conventional comparators, of which comparator 116 is representative, and a conventional decoder 118 that produces a multi-bit digital output 120. As is well known in the art, input voltage 110 is applied simultaneously to each comparator 116. In addition, fractional portions of the reference voltage 112 are applied to the comparators 116 by dividing the reference voltage 112 in equal increments (or thresholds) by the resistors 114. The output of each comparator 116 is applied to the decoder 118 which decodes such received inputs into a multi-bit digital output 120 representative of the input voltage 110. Although a single-ended structure is shown in FIG. 1 and throughout this discussion, in practice a fully differential structure can be used. ADCs for operation at high frequencies, however, require a large amount of integrated circuit area and have high power consumption, and all such requirements increase as the number of bit of resolution of the ADC increases. For example, a 6-bit full flash ADC requires about 26=64 voltage comparators. In a CMOS implementation of a full flash ADC, these comparators are normally implemented using conventional auto-zero voltage comparators. An auto-zero voltage comparator, however, requires a two-phase clock for auto-zeroing in the first phase, and for actual signal comparison in the second phase. Unfortunately, such two-phase design limits the maximum achievable operating frequency to a factor of two lower than otherwise possible, other factors being equal, if non-auto zero voltage comparators are employed. Non-auto zero voltage comparators, such as those used in full flash ADCs implemented in Bipolar or BiCMOS integrated circuit processes, are generally not practical for implementation in standard CMOS processes because device mismatches (e.g., input offset voltage) of CMOS voltage comparators tend to be much higher than for Bipolar equivalents. CMOS voltage comparators with low input offset voltage can usually only be obtained using complex circuitry that requires large integrated circuit area with associated higher power consumption, and generally lower conversion speed. Therefore, it is desirable to provide a high resolution ADC that has small die size and low power consumption, and that avoids the effects of operational mismatches. Accordingly, the full flash ADC of the present invention includes a plurality of comparators and a referencing scheme that effectively cancels out the input offset voltages of the comparators. The input offset voltage of each of the plurality of comparators is obtained by performing a self calibration process on the ADC during, for example, power up. Then, the input offset voltage for each of the comparators is stored in a look-up table. When the ADC is used, the look-up table provides offset correction to the normal reference voltages for each comparator. In one embodiment of the present invention, the offset look-up table controls a digital-to-analog converter (xe2x80x9cDACxe2x80x9d). In addition, a track/hold (xe2x80x9cT/Hxe2x80x9d), circuit, also known as a sample/hold circuit, is connected to a reference input of each comparator. The T/H circuit receives its input from the DAC and holds received voltages for application to its associated comparator as a first or reference input. Each comparator receives the analog input voltage as its second input and the outputs of the comparators are supplied a conventional decoder. The look up table, in combination with a T/H controller and the DAC, operates each T/H circuit to provide a voltage equivalent to the reference voltage corrected by the input offset voltage of the associated comparator. After the correct reference voltages are loaded into the T/H circuits, the analog input signal is applied to all of the comparators. Each comparator produces an output signal indicating whether the magnitude of the input signal is, for example, greater than the magnitude of the corrected reference voltage. The decoder receives such comparator outputs and decodes the outputs into a representative multi-bit digital output signal.
{ "pile_set_name": "USPTO Backgrounds" }
One of the most common and life-threatening heart irregularities is ventricular fibrillation in which the heart is unable to pump a significant volume of blood. When such an irregularity occurs, serious brain damage and death will invariably result unless a normal heart rhythm can be restored within a few minutes. Ventricular fibrillation can occur as a result of a heart attack but may also be caused by accidental electric shock or due to severe stress, such as may accompany surgery, drowning or the like. The most effective treatment for restoring a normal rhythm to a heart muscle experiencing ventricular fibrillation is the application of a strong electric shock to the victim using a cardiac defibrillator. Cardiac defibrillators are medical devices for producing and delivering such shocks and have been successfully used for many years. Conventional external cardiac defibrillators accumulate an electric charge on a storage capacitor and, when a switching mechanism is closed, transfer the stored energy in the form of a large current pulse to a patient. The switching mechanisms used in most defibrillators comprise heavy-duty electro-mechanical relays. Typically, the relays are responsive to a discharge control signal that actuates the relay to complete an electrical circuit between the storage capacitor, a waveshaping network, and a pair of defibrillation electrodes attached to the patient. The cardiac defibrillation pulse, which is delivered to the patient by an energy transfer circuit that includes a electro-mechanical relay, generally comprises a single pulse having a damped sinusoidal shape that starts when the relay closes. Alternately, a cardiac defibrillation pulse may have an exponential shape that starts when the relay contact closes and stops when the relay contact opens. However, it has been found that other types of defibrillation currents may be more effective in terminating ventricular fibrillation. For example, it may be beneficial to apply a series of cardiac defibrillation current pulses to the patient in rapid succession, or it may be beneficial to maintain the magnitude of the defibrillation current flowing through the patient at a specific level. The use of electro-mechanical relays in an energy transfer circuit do not allow such types of defibrillation currents to be applied to the patient because the relay contacts cannot be switched closed and open fast enough. This slow response time does not allow the use of feedback signals obtained from the patient to be used in controlling the amount of energy delivered. Therefore, what is needed is an energy transfer circuit that can replace the electro-mechanical relays with a fast acting solid state switching device. The problem with using solid-state switching devices in a cardiac defibrillation circuit is that such devices allow a certain amount of leakage current to flow even when the devices are in a nonconducting state. Heretofore, safety reasons have made it impractical to use solid-state circuitry devices in a cardiac defibrillator circuit without the additional isolation of electro-mechanical relays. Therefore, it is desirable to provide an energy transfer circuit that employs a solid state switch in combination with a current shunt which can divert the leakage current away from the patient electrodes until a defibrillation pulse is delivered to the patient. It is also desirable to provide an energy transfer circuit including a feedback control to regulate the delivery of a cardiac defibrillation pulse to the patient. Finally, it is desirable to provide an energy transfer circuit that can deliver predetermined defibrillation current waveforms to the patient.
{ "pile_set_name": "USPTO Backgrounds" }
The present invention relates to a method of and an apparatus for assembling an electronic component, such as a ladder-type filter, having a plurality of inner parts being held under pressure and stored in a case. FIGS. 1 and 2 show a conventional four-element ladder-type filter. This filter comprises two series ceramic resonators 2 and 3, two parallel ceramic resonators 4 and 5, an input terminal plate 6, a ground terminal plate 7, an output terminal plate 8, an internal connection terminal plate 9 and a plate spring 10, which are stored in a box-type case 1. Outlines of these inner parts are substantially equivalent to each other. FIG. 3 is a circuit diagram showing the electrical connection thereof. The input terminal plate 6 is provided on its one major surface with a protruding portion 6a which is in pressure contact with a central portion of the series resonator 2. The ground terminal plate 7 is provided on both its major surfaces with protruding portions 7a and 7b which are in pressure contact with central portions of the two parallel resonators 4 and 5 respectively. The output terminal plate 8 is provided on both its major surfaces with protruding portions 8a and 8b which are in pressure contact with central portions of the parallel resonator 5 and series resonator 3 respectively. The internal connection terminal plate 9 has a portion 9a which is inserted between the series resonator 2 and parallel resonator 4, a portion 9b which is arranged next to an outer major surface of the series resonator 3, and a coupling portion 9c which connects these portions 9a and 9b. Since the plate spring 10 is interposed between the portion 9b and an inner side surface of the case 1, the terminals 6-9 and the resonators 2-5 are in pressure contact with each other to be electrically connected. A cover sheet 11 is inserted in an opening of the case 1, and a cavity which is defined by the cover sheet 11 and the opening is filled up with filler 12 such as resin, whereby the opening of the case 1 is sealed. Lead portions 6b, 7c and 8c of the input terminal 6, the ground terminal 7 and the output terminal 8 protrude from the opening of the case 1. In the conventional ladder-type filter having the aforementioned structure, the resonators 2-5 utilizing surface-directional spreading vibration mode are merely in point contact with the terminals 6-9 at the central protruding portions 6a, 7a, 7b, 8a and 8b. Therefore, it is necessary to manually insert the resonators 2-5 and the terminals 6-9 into the case 1 one by one with tweezers since these parts easily get slanted due to the point contact. As a result, working efficiency of the assembling operation becomes quite low. Further, since it is necessary to bring all parts into close contact with each other in the final stage of assembly, much time is required for inserting the final part into a very small clearance. Moreover, electrode surfaces of the resonators 2-5 may be damaged by the assembling operation. Thus, the conventional method of assembling electronic component such as a ladder-type filter had a problem of low working efficiency caused by manual assembly.
{ "pile_set_name": "USPTO Backgrounds" }
The invention relates to a lubricant or hydraulic pump for delivering oil or the like with a housing, an inlet provided on or in the housing, a drive element mounted rotatably in the housing, and at least one pump element, which interacts with the drive element for delivering oil or the like. Such pumps are used for delivering lubricant to lubricating points of machines or for delivering a hydraulic fluid. For example, such pumps are used for presses, in particular, packing presses. The configuration of such pumps is complicated, in part, by many different components. This makes the production of such pumps more expensive. In addition, it is often not possible to quickly adapt the number of pump elements in use to the individual requirements of the application, so that the possible applications of known pumps are somewhat limited. The task of the present invention is therefore to create a lubricant or hydraulic pump of the type named above, which can be produced in an especially economical way and which can have various uses.
{ "pile_set_name": "USPTO Backgrounds" }
Field of the Invention The present invention relates to air handling systems, such as heating and cooling systems and, more particularly, to a gravity exhaust air shutter. Background Information Many commercial buildings are heated and/or cooled by one or more self-contained, packaged air-conditioning units. Many of these units are installed on the roof tops of such buildings. These units typically include cooling equipment, air handling fans, and may further include dehumidifiers and heating equipment. The units receive return air from the building, cool (or heat) the received return air, and supply the cooled (or heated) return air back to the building. Some units also include exhaust fans for discharging some portion of the return air from the building. The exhaust fans are typically mounted next to an exhaust opening of the unit. In a prior design, a series of is louvered slats are disposed at the opening. The louvered slats may extend the width of the opening, and may overlap each other when in the lowered, e.g., closed, position. Pins on the ends of the louvered slats may be mounted in holes formed in the sides of the exhaust opening, thereby allowing the louvered slats to swing open when the exhaust fans are turned on. Outside air often leaks into the unit through these louvered slats. The occurrence of outside air leaking into the unit can degrade its performance. Accordingly, a need exists to improve the performance of packaged units and other air handling systems.
{ "pile_set_name": "USPTO Backgrounds" }
As shown in U.S. Pat. No. 4,113,268 dated Sept. 12, 1978, a sealing assembly for a butterfly valve is illustrated in which a soft primary seal is positioned between a pair of spaced metal secondary seals. The secondary metal seals are held out of contact with the seating surface of the butterfly disc by the soft seal under normal operating conditions. Only when the soft primary seal is destroyed do the secondary metal seals engage the disc in sealing relation. A separate resilient backing ring urges the primary soft seal and the secondary metal seals inwardly toward the seating surface of the disc. U.S. Pat. No. 4,162,782, issued July 31, 1979 to Ronald D. Wilkins and filed Apr. 10, 1978 for "Seal Assembly for Butterfly Valve" discloses a seal assembly mounted in an annular groove about the flow passage. The seal assembly includes a metallic body having a pair of outer legs seated in the bottom of the groove and a pair of inner legs contacting the outer periphery of the valve disc. A soft seal is positioned between the inner legs and the inner ends of the inner legs contact the valve disc to provide metal sealing surfaces. Upon contact with the disc, the inner legs are urged radially outwardly. However, the inner legs do not have a large radial movement and a substantial frictional force is provided by contact with the inner legs upon radially outward movement which increases the operating torque for moving the valve between open and closed positions. With certain types of metal finishes or coatings on the valve disc, some scratching of the sealing surface on the valve disc might occur with a high frictional contact, particularly if the seal has a so-called knife edge metal contact surface.
{ "pile_set_name": "USPTO Backgrounds" }
Rotaviruses are consistently shown to be the single most important cause of severe diarrhea of infants and young children in both developed and developing countries. The consequences of rotavirus diarrhea are staggering as they account for up to 592,000 deaths annually in the under 5-year age group, predominantly in the developing countries (Parashar et al, Emerg. Infect. Dis., 2003, 9:565-572). It has recently been estimated that 1 in 200 children in developing countries will die from rotavirus diarrhea (Glass et al, Lancet, 2004, 363:1547-1550). In the United States, in the under 5-year age group, it was estimated that annually rotaviruses are responsible for 2,730,000 episodes of diarrheal illness, 410,000 visits, to a physician, 160,000 emergency room visits, 50,000 hospitalizations, and 20 deaths (Tucker et al., JAMA, 1998, 279:1371-1376). Thus, the need for a rotavirus vaccine in both developed and developing countries has received national and international endorsement. WO2007009081 had proposed the use of a hexavalent bovine rotavirus (UK)-based vaccine for developing countries to cover not only the standard serotypes G1 through G4 but also emerging serotypes G8 and G9. Development of Hexavalent Rotavirus vaccine comprising of standard G1-G4 strains and G9, G8 strains for a broader degree of protection has been previously discussed see Albert Z. Kapikian et. al. “A Hexavalent human Rotavirus-bovine Rotavirus reassortant vaccine designed for use in developing countries”; National Institutes of Health; Journal of infectious diseases; 2005; 192; S22-9 Rotavirus strains may lose viability during drying process and storage. It has been reported previously that lyophilization causes upto 30% loss in virus potency. Protein formulations containing sucrose-glycine combination have been described by Wei Liu et. al. “Freeze drying of Proteins from a sucrose-glycine excepient system: Effect of formulation composition on initial recovery of protein activity”; AAPS Pharm Scitech; Feb. 11, 2005; 6(2); E150-E157 A lyophilized rotavirus vaccine formulation according to U.S. Pat. No. 6,616,931, U.S. Pat. No. 6,403,098 & U.S. Pat. No. 5,932,223 comprising a strain of rotavirus about 1×105 to about 1000×105 pfu/mL, a sugar 1-20% (w/v), Phosphate about 0.05-2 M and a) 0.5%-1.25% of recombinant human serum albumin or b) 0.1%-1.25% of an amino acid (glutamate, glutamine or arginine). Further the patents also discuss improved stability of lyophilized rotavirus formulation by including glycine (1%) in the sucrose/mannitol stabilizer. There remains a distinct need for rotavirus vaccine formulations with improved viability and stability. None of the prior art stabilizers improve viability & stability. Further for worldwide distribution of rotavirus vaccines, it is necessary to formulate vaccines such that they are stable under a variety of environmental conditions. Surprisingly, it has now been discovered that a) stabilizer composition comprising of combination of Sucrose and Glycine and b) optimal lyophilization cycle results in a Rotavirus formulation with a moisture content less than 3% and 100% individual virus preservation.
{ "pile_set_name": "USPTO Backgrounds" }
Synchronization of a decoding and presentation process for received bitstreams is a particularly important aspect of real-time digital data delivery systems. For example, the Moving Pictures Experts Group (MPEG) has promulgated several standards relating to digital data delivery systems. The first, known as MPEG-1 refers to ISO/IEC standards 11172-1 (Systems), 11172-2 (Video), 11172-3 (Audio), 11172-4 (Compliance Testing), and 11172-6 (Technical Report). The second, known as MPEG-2, refers to ISO/IEC standards 13818-1 (Systems), 13818-2 (Video), 13818-3 (Audio), 13818-4 (Compliance). The MPEG standards address the timing and synchronization issues for decoders of MPEG data streams (e.g., video, audio, data, and the like) as follows. A sample of a 27 MHz clock is transmitted in a program clock reference (PCR) field of a transport stream packet. The PCR indicates a time when the transport decoder is expected to have completed reading the PCR field. The phase of the local clock running at the decoder is compared to the PCR value in the bit stream at the instant at which it is obtained to determine whether the decoding process is synchronized. In general, the PCR from the transport bit stream does not directly change the phase of the system clock of the decoder, but only serves as an input to adjust the clock rate of a Voltage Controlled Crystal Oscillator (VCXO). This adjustment is accomplished by comparing the PCR sample to the output frequency of the VCXO and responsively adjusting the VCXO until the two frequencies match. The system clock of the decoder is then derived from the VCXO. Thus, the system clock of the decoder is adjusted and locked to the 27 MHz clock of the encoder or transmitter via the transmitter clock sample included in the PCR. A disadvantage of the prior art decoder arrangement is the relatively high cost of the VCXO. Therefore, a need exists in the art for a method and apparatus which is able to perform timing functions in a data delivery system using an inexpensive oscillator, e.g., a substantially fixed frequency oscillator.
{ "pile_set_name": "USPTO Backgrounds" }
The invention relates to the problem of optimizing the transport of physical objects. The International patent application with the International publication number WO 03/035282 A2 and the corresponding European patent EP 1 455 959, entitled to the Deutsche Post AG, describe a method for processing objects wherein information located on at least one surface of the object is detected. The processing of the objects is characterized in that address information determined by means of the information located on the surface of the object is compared with available address information in a databank or in a database drawn up there from. The United States patent application US 2005/0021389 A1 also describes a method and a system for calculating an environmental score for a business unit. It also includes a computer system for calculating a score for separately accountable business units, the score being indicative of a level of unaccounted aspects of external environmental cost of economic activities. The United States patent publication US 2005/0052810 A1 describes a method and an apparatus for calculating an environmental indicator and recording medium with calculation program recorded there on. According to this document database and data table are stored in a memory. The database has data associated with the part lists and product specifications of products in conjunction with products identification codes, whereas the data table has processing yields and environmental indicator factors in conjunction with material codes which respectively indicate the material of each part constituting a product. Afterwards processing yield and environmental indicator factor for every material code are calculated by referring the data table, the material codes relating to the parts corresponding to the part numbers which have been extracted.
{ "pile_set_name": "USPTO Backgrounds" }
1. Field of the Invention The present invention relates to devices for monitoring around a vehicle in which target extraction is carried out by subjecting an image captured by an infrared camera device to a binary thresholding method. 2. Description of Related Art Devices for monitoring around a vehicle have been proposed in which objects that may collide with the vehicle, such as pedestrians, are extracted from a picture around the vehicle captured by an infrared camera, and such information is supplied to a driver of the vehicle. In these devices, the probability that the vehicle will hit an object, such as a pedestrian, is calculated based on the relative distance between the vehicle and the object, and the relative speed of the vehicle. An example of such devices for monitoring around a vehicle which extract an object, which may collide with the vehicle, from an image around the vehicle captured by an infrared camera is as follows. That is, in the device, the captured infrared image is subjected to a binary (2-level) thresholding process, and an area to which bright (white) portions are concentrated is searched for. Then, it is determined whether the area is a head portion of a pedestrian by using the aspect ratio (ratio of length to width) and the sufficiency rate of the area, and further calculating the distance between the vehicle and the area using the actual surface area and the position of the center of gravity in the image. If the area of the head portion of a pedestrian is determined, an area which forms the body of the pedestrian is determined by calculating the height of the pedestrian in the image based on the distance between the area determined to be the head portion and the camera, and an average height of an adult. These areas are displayed to be distinguished from the other regions of the image. In this manner, the position of the entire body of the pedestrian in the infrared image is determined, and this information is displayed for the driver so as to effectively assist the vision of the driver (refer to Japanese Unexamined Patent Application, First Publication No. Hei 11-328364, for example). However, using the thresholding method, only the head portion, a part of the head portion, or the entire or only an upper half or a lower half of the body, of a pedestrian may be extracted from an infrared image, depending on the effects of a hat or clothing the pedestrian is wearing, or of the environment surrounding the pedestrian, and thus the shape of the pedestrian obtained by the binary thresholding method becomes ambiguous. Also, when a vehicle is running, due to the influence in change in the shape of a road ahead, or the pitting of the vehicle, the height of a pedestrian, from a child to an adult, is generally detected to be different from his/her real height. Accordingly, since the barycentric coordinates of targeted objects, such as pedestrians, in the picture cannot be fixed with respect to the distance. Therefore, it is not possible to extract the targeted objects stably such as pedestrians which may collide with the vehicle, if the extraction is carried out based on the shape of at least the pedestrians' height, head, body as in the above-mentioned conventional device.
{ "pile_set_name": "USPTO Backgrounds" }
Search engines generally attempt to provide users with fast, accurate, and timely search results. Several search engines exist with varying interfaces and algorithms that assist computer users in finding resources stored on one or more computing systems, including a network of computing systems. With the advent of the Internet, search engines have grown in speed and functionality, to accommodate demand for finding the billions of resources stored in computing systems connected across the Internet. Web search engines, for instance, are search engines designed to search for information on the Internet. Typically, a user submits a search query specifying certain keywords, criteria, or conditions and the search engine consults one or more indexes to determine which resources, known to the search engine, likely satisfy the search query. Results of the search, also known as “hits,” can be returned to the user. In some cases, the user can access or request the resources included in a listing of hits directly from the listing, for instance, through the selection of a corresponding hyperlink. In some existing search engines, multiple different users independently submit search queries and interact with search results generated in response to the queries. Users' queries and interactions with the search results can be collected and stored as histories and can be associated with particular users or categories of users. Some modern search engines attempt to take an individual user's history or profile into account when returning results of a given search query. For instance, a geographic location associated with the user can be employed to return, if possible, search results that accord with that location. Historical trends observed for a given user can also be considered in a search algorithm as well as the aggregate historical trends of multiple independent users to assist in generating search result sets more responsive to particular queries. For instance, in some modern search engines, intelligence gleaned from previous searches performed by one or more users can be used to influence both the search criteria recommended and the search results returned for another user in a subsequent search. Like reference numbers and designations in the various drawings indicate like elements.
{ "pile_set_name": "USPTO Backgrounds" }
A Schmoo plot, also known as a Shmoo plot, schmoo, or shmoo, is a technique that is used to characterize the results of integrated circuit testing. Measurements performed on a device under test (DUT) may be performed as one or more parameters are varied. If a first parameter P1 is tested over a first range of M values and a second parameter P2 is tested over a second range of N values, then there are M*N possible measurements that can be made. The results of testing the DUT at these M*N values of P1 and P2 may be plotted in a two-dimensional grid. The resulting plot is called a Schmoo plot, and indicates regions of successful and unsuccessful testing. Schmoo plots can have more than two dimensions, such as when a DUT is tested as three parameters are varied over a 3-D grid to produce a 3-D Schmoo plot. Determining the Schmoo plot can become very time-consuming in a finely grained mesh since each measurement point must be computed sequentially. The time-consuming nature of Schmoo plots is problematic, and can result in failure analysis of a small number of DUTs.
{ "pile_set_name": "USPTO Backgrounds" }
1. Technical Field This disclosure relates to a constant voltage circuit, and more particularly to a constant voltage circuit capable of making a quick response to a wide range of output currents such as a minute current and a large current, and capable of stable operation with high efficiency. 2. Description of the Related Art In electronic devices such as portable phones, mobile PCs, and car navigation systems, a constant voltage power source having a constant voltage circuit and capable of supplying a stable voltage is used as a power source. When using such a constant voltage power source in a device with a large output current, the constant voltage power source is required to have a circuit configured to realize a high speed response by improving a ripple removing ratio and a load transient response. For example, when the constant voltage power source is used in a device with a wide range of output current, such as a portable phone having an operation mode and a standby mode, a circuit configuration capable of receiving a maximum output current is required. As a result, a current consumption is increased as a whole. In the standby mode of the portable phone, in which a high ripple removing ratio and a high load transient response are not required, an unnecessary current is consumed, which results in increasing the wasted current. In view of this, a constant voltage circuit for suppressing this wasted power consumption has been suggested. Each of Patent Documents 1 and 2 discloses a constant voltage circuit configured to increase or decrease a bias current supplied to a differential amplifier in the constant voltage circuit depending on the amount of output current. FIG. 8 shows the constant voltage circuit disclosed in Patent Document 1. In FIG. 8, a constant voltage circuit 101 includes a reference voltage circuit Vref, a differential amplifier circuit 102, a bias current generating circuit 103, and an output circuit 104. In this circuit, a PMOS transistor M7 and an output transistor M1 form a current mirror circuit. Therefore, a drain current in proportion to a drain current (output current) of the output transistor M1 is generated in the PMOS transistor M7. This current is supplied as a drain current of an NMOS transistor M8. Since the NMOS transistor M8 and an NMOS transistor M9 form a current mirror circuit, a drain current of the NMOS transistor M9 is in proportion to the drain current of the output transistor M1. The drain current of the NMOS transistor M9 is a part of a bias current of the differential amplifier circuit 102, therefore, the bias current of the differential amplifier circuit 102 increases and decreases in accordance with an increase and a decrease of the output current. In this manner, the bias current of the differential amplifier circuit 102 is increased and decreased in accordance with the increase and decrease of the output current. Therefore, a response speed is increased when the output current is increased. In this manner, the current consumption and the response speed are set appropriately. [Patent Document 1] Japanese Patent Application Publication No. 3-158912 [Patent Document 2] Japanese Patent Application Publication No. 2006-99526 In the constant voltage circuits configured to change the bias current of the differential amplifier circuit in accordance with the output current as disclosed in Patent Documents 1 and 2, an operation of the constant voltage circuit becomes unstable when the output current is small. That is, for example, a constant voltage power source having a large output transistor a capable of outputting an output current of 1 A or more can be stably operated when the output current is large. However, this constant voltage power source cannot be stably operated when the output current is small since a bias current of a differential amplifier circuit becomes small and a phase margin is decreased. Moreover, there is a problem in that a response speed is extremely low when the bias current is small. This is because a transistor having a large ratio of gate width to gate length and thus having large gate capacitance is used as an output transistor to realize an operation with a large current. When a bias current is small, it takes time to charge and discharge the gate capacitance. Therefore, the response speed is drastically decreased when the output current is small.
{ "pile_set_name": "USPTO Backgrounds" }
Increasing the speed of operation of electronic devices, such as computer systems, requires increasing the speed of not only processors, digital logic and data storage components, but also the speed of buses used to couple such components together. However, increasing the speed of operation of a bus such that data is transferred across that bus more quickly often entails the use of combinations of transmitter and receiver circuit designs that consume more power. In some cases, this increased power consumption is the result of having to use components in transmitter and/or receiver circuits that have the desirable characteristic of being operated at higher speeds, but which require more power to switch so quickly between states needed to transfer a binary value of 0 or 1. In other cases, this increased power consumption is the result of having to couple together components that have the higher speed characteristic, but which interact with other needed components in ways that may create other undesirable conditions that must be overcome through the use of more power, such as higher parasitic capacitive loads. In still other cases, this increased power consumption is the result of combinations of transmitter and receiver design or referencing to voltage levels such that the transmitters and receivers are required to use the same voltage level as a power rail, even though the core voltage level at which one of the devices internally operates is a lesser voltage. The fact that the speed of processor, support logic and data storage components have been increased while being necessarily based on differing semiconductor design and/or process technologies has also created voltage level incompatibility issues with processors and digital logic components typically being designed to operate at ever lower voltage levels (currently 1 volt or less), while data storage components, especially dynamic random access memory (DRAM) devices typically operate at higher voltages (currently 1.3 volts or higher). This mismatch in voltages arising from differing semiconductor process technologies typically results in processors and/or digital logic components having to employ transmitter circuit designs that not only transmit data across a bus to data storage components, but which also internally convert from the lower core voltage within a processor or digital logic component to a higher external voltage that matches the I/O voltage employed by a data storage component, because prior art transmitter and receiver circuits often do not work correctly unless both the transmitter and receiver circuits employ the same voltages, themselves. To support this conversion between voltages, a processor or digital logic component must be coupled to two different voltage rails, one for the core and the other for transmitters and/or receivers, must either use special high voltage tolerant transistors to handle the higher voltage within the lower voltage silicon technology which increases process technology costs through added process complexity, or use cascode transistor techniques to handle the higher voltage which increases costs through taking up more space on a silicon die. The detrimental effects of higher power consumption and higher silicon technology costs arising from such approaches where both transmitters and receivers are required to operate at the same higher voltage are incurred regardless of whether the signaling across a bus between transmitter and receiver circuits is entirely single-ended, entirely differential, or a mix of the single-ended and differential. These difficulties with voltage level incompatibility in current practice are illustrated by FIGS. 1a, 1b and 1c depicting prior art transmitter and receiver circuit designs. It should be noted that although for the sake of simplicity of discussion, FIGS. 1a-c depict only unidirectional configurations, these same issues arise and apply to bidirectional configurations, as well. In FIG. 1a, transmitting device 120 (such as a memory controller IC) employs multiple ones of single-ended transmitter 130 (although only one is shown for sake of clarity) to transmit addresses, commands and/or data across bus 150 to single-ended receiver 170 employed by receiving device 160 (such as a memory IC). Switches 131 and 139 receive data from other portions of transmitting device 120 and drive either a high or low voltage level onto a signal line of bus 150 through resistors 132 and 136, respectively, while resistors 172 and 176 are employed by receiving device 160 to terminate this same signal line of bus 150, referencing receiver Vcc and ground, respectively, at or near the point at which this same signal line is coupled to the input of single-ended receiver 170. Capacitors 134 and 137 are parasitic capacitors, i.e., capacitive loads arising from the connection of switches 131 and 139 to this same signal line of bus 150, thereby slowing down the speed at which the signal state of this signal line may be changed in transmitting data. The presence of resistors 132 and 136 does mitigate this undesirable effect on this signal line, but only to a limited extent, since mitigation to a greater degree would require a higher resistance value which would, in turn, defeat the ability of switches 131 and 139 to drive this signal line. Also, it is common for single-ended transmitter 130 to be designed to conform to a specification of electrical characteristics for a signal line, including signal line 150, such that the resistance of resistors 132 and 136 is often dictated by such a specification, and therefore, cannot be changed. Another undesirable effect of this configuration of transmitter and receiver design is that transmitter Vcc and receiver Vcc must be of the same voltage level for high and low values that distinguish between binary 1 and 0 values to be correctly detected by single-ended receiver 170. FIG. 1b depicts a somewhat different design for single-ended transmission of data from FIG. 1a, but despite the design differences, largely the same previously discussed problems are presented again. Capacitor 137 again depicts a parasitic capacitor on a signal line of bus 150, again slowing the speed at which the state of that signal line may be changed, and again, the configuration of transmitter and receiver design requires that transmitter Vcc, the receiver Vcc and the power rail to which termination may be coupled must all be of the same voltage level. The same difficulties would continue to exist even if transistor 139 were coupled to a transmitter Vcc and resistor 172 were coupled to ground. Furthermore, despite the depiction in FIG. 1c of a differential receiver in contrast to the use of single-ended configurations in FIGS. 1a and 1b, the same difficulty of differential receiver 180 needing to be supplied with a receiver Vcc that matches the same voltage level as supplied to whatever differential transmitter may drive the pair of signal lines received by differential receiver 180 from across bus 150 still exists.
{ "pile_set_name": "USPTO Backgrounds" }
The secretion of hypophysial anterior lobe hormone is regulated by the peripheral hormone secreted by each target organ and the secretion-promoting or secretion-suppressing hormone secreted by the hypothalamus, which is the center superior to the hypophysial anterior lobe, and this group of hormones hereinafter generically referred to as hypothalamic hormone in this specification. To date, nine hypothalamic hormones have been identified, for example, thyroid-stimulating hormone-releasing hormone (TRH), and gonadotropin releasing hormone [GnRH, also known as luteinizing hormone releasing hormone (LH-RH)], etc. It is conjectured that these hypothalamic hormones exhibit their hormone actions etc. via receptors assumed to be present in the hypophysial anterior lobe, and analyses of receptor genes specific to these hormones, including humans, are ongoing. Antagonists or agonists that act specifically and selectively on these receptors would therefore regulate the action of hypothalamic hormones and hence regulate the secretion of hypophysial anterior lobe hormone. As a result, such antagonists or agonists are expected to prevent or treat diseases depending on these hypophysial anterior lobe hormone. Known compounds possessing GnRH-antagonizing activity include GnRH-derived linear peptides (U.S. Pat. Nos. 5,140,009 and 5,171,835), a cyclic hexapeptide derivative (JP-A-61-191698), a bicyclic peptide derivative [Journal of Medicinal Chemistry, Vol. 36, pp. 3265-3273 (1993)], and so forth. Non-peptide compounds possessing GnRH-antagonizing activity include compounds described in WO 95/28405, WO 97/14697 and WO 97/14682, etc. ZA 86/9289 describes 1,4-dihydro-1-ethyl-7-phenylpyrrolo[1,2-a]pyrimidin-4-one in Example 22. Peptide compounds pose a large number of problems to be resolved with respect to oral absorbability, dosage form, dose volume, drug stability, sustained action, metabolic stability etc. There is strong demand for an oral GnRH antagonist, especially one based on a non-peptide compound, that has excellent therapeutic effect on hormone-dependent cancers, e.g., prostatic cancer, endometriosis, precocious puberty etc., and that does not show transient hypophysial-gonadotropic action (acute action).
{ "pile_set_name": "USPTO Backgrounds" }
Economics of distribution of agricultural chemicals, such as herbicides in general and glyphosate formulations in particular, can be much improved through provision of a high “loading” of active ingredient in the formulation, that is, the amount of active ingredient that can be accommodated in a container of given capacity. Glyphosate is an acid that is relatively insoluble in water (1.16% by weight at 25° C.). For this reason it is typically formulated as a water-soluble salt in aqueous solution. A useful alternative is to prepare glyphosate as a dry salt in powder or granular form. For example, a dry water-soluble granular formulation of glyphosate ammonium salt can have a glyphosate acid equivalent (a.e.) content as high as about 86% by weight. This would appear at first sight to provide an excellent solution to the challenge of packing more glyphosate into a container of given capacity. Unfortunately the benefit of a dry glyphosate formulation in this regard is more limited than one might expect, because such a formulation tends to have low bulk density. Also, many end-users and many distributors prefer a liquid product because of flexibility in handling, thus a need remains for high-loaded liquid formulations of glyphosate. U.S. Pat. No. 6,544,930 to Wright discloses an approach to meeting this challenge. According to this approach, a concentrated aqueous solution of glyphosate, predominantly in the form of one or a mixture of the potassium and monoethanolammonium (MEA) salts thereof, was provided, it having been determined that such a solution had an unexpectedly high specific gravity, permitting more glyphosate a.e. to be delivered in a container of given capacity than was previously attainable using the isopropylammonium (IPA) salt in widespread commercial use, for example as Roundup® herbicide of Monsanto. Unfortunately, glyphosate potassium salt, especially when formulated at high concentration in aqueous solution, brings some challenges of its own. For example, where (as often) it is desired to coformulate a surfactant with the glyphosate, physical incompatibility of the surfactant with the glyphosate salt can limit the options available. Whereas a wide range of surfactant types are compatible with glyphosate IPA salt, only a few types have been found to be compatible with glyphosate potassium salt, in particular where the salt is present at high concentration. See above-cited U.S. Pat. No. 6,544,930, col. 9, lines 6-13. International Patent Publication No. WO 01/26469 discloses that aqueous formulations of glyphosate, including highly concentrated formulations, can be prepared using a mixture of glyphosate IPA and ammonium salts at a weight ratio (expressed on a glyphosate a.e. basis) of 80:20 to 97:3. Such formulations are said to exhibit reduced viscosity, leading to greater ease of pumping and handling. International Patent Publication No. WO 03/013241 proposes, inter alia, a glyphosate composition comprising IPA and potassium cations in a mole ratio of 1:10 to 30:1, “more preferably less than 15:1 and greater than 1:2”, reportedly as a means to improve bioefficacy over compositions of a single glyphosate salt. Publications cited above are incorporated herein by reference. Considering the variety of conditions and special situations under which glyphosate herbicides are used around the world, there remains a need for aqueous concentrate formulations of glyphosate, including surfactant-containing formulations, providing benefits under at least some of those conditions and situations. There is an especial need for such formulations having high glyphosate loading, for example at least about 400 g a.e./l.
{ "pile_set_name": "USPTO Backgrounds" }
EP-A-0 51 351 disclose an Mg, Al mixed (hydr) oxide catalyst having an Mg:Al molar ratio above 3 and preferably in the range from 3-10. The article of H. Schaper et al. in Applied Catalysis, 54, (1989) 79-90, discloses the same catalyst. This catalyst has a hydrotalcite structure, consisting of brucite type layers in which part of the bivalent ions (Mg) are replaced by trivalent ions, alternated by interlayers which contain water and anions to compensate for the excess charge of the trivalent ions. The preparation of such catalysts is disclosed. Due to the basic properties such catalysts are considered of special interest for base-catalyzed reactions, such as polymerization of propylene oxide, double-bond isomerizations of olefins such as 1-pentene, and aldol condensations. Exemplified is double-bond isomerization of 1-pentene using an Mg, Al mixed oxide catalyst having an Mg:Al molar ratio of 5 and 10. At increasing molar ratio the conversion rate decreases. The article of Watanabe, Y. et al. in Microporous and Mesoporous Materials 22 (1998) 399-407, discloses the use of Mg—Al hydrotalcite catalysts having a molar ratio of 1.8-2.5 for the methanolysis of ethylene carbonate for the production of dimethyl carbonate. EP-A-0,478,073 describes a process for preparing a dialkyl carbonate which comprises contacting an alkylene carbonate with an alkanol in the presence of a mixed metal oxide catalyst or a modified bimetallic or polymetallic catalyst under conditions effective to produce the dialkyl carbonate. In the examples, a magnesium/aluminium mixed metal oxide catalyst having a Mg:Al ratio of 3:1 was employed. In JP-A-06/238165, a process is described wherein an alkylene carbonate and an alcohol are subjected to transesterification in presence of a catalyst to produce a dialkyl carbonate. A combination of Magnesium oxide and another metal oxide other than magnesium was used as catalyst in an atomic ratio in the range of 1000:1 to 20:1 of magnesium to the other metal. The present invention has for its object to provide a method for the catalytic conversion of alkylene carbonate having an improved conversion rate and improved yield, while having limited leaching of metal from the catalyst.
{ "pile_set_name": "USPTO Backgrounds" }
1. Field of the Invention The present invention relates to a process and an evaporation apparatus, and more particularly, to an evaporation apparatus capable of preventing a sag phenomenon of a substrate. 2. Discussion of Related Art A contemporary evaporation apparatus may include a substrate supporting unit, and the substrate supporting apparatus includes a chamber and a substrate supporter for supporting two sides of a substrate in the chamber, when the substrate enters the chamber. The substrate supporter may move repeatedly up and down driven by a lever assembly provided at the exterior of the chamber. Referring to the method of supporting a substrate using the substrate supporter, the lever assembly is used to allow the substrate supporter to move down before the substrate enters the chamber. When the substrate enters the chamber, the substrate is loaded onto the substrate supporter. Then, the substrate supporter moves upwards. More particularly, the substrate supporter supports the substrate by holding one side and another opposite side facing the one side of the substrate with the substrate supporter physically contacted with a bottom surface of the substrate, the bottom surface facing a lower surface of the chamber. In other words, two opposite sides of the substrate are supported by the substrate supporter with the substrate being loaded onto the substrate supporter. When the substrate supporter supports two opposite sides of the substrate as described above, however, it is impossible to support a large-area substrate uniformly across the width of the substrate. That is to say, a central region of the substrate may sag due to the increased load of the substrate when the substrate is manufactured on a large scale. In other words, a large size substrate may become deformed due to the earth's gravity.
{ "pile_set_name": "USPTO Backgrounds" }
1. Field The present inventive concept relates to a heat radiation sheet and a method of manufacturing the same. 2. Description of the Related Art Heat may he generated in an electronic device by electronic parts such as wirings, terminals, and chips. Heat generated by electronic parts may shorten the lifetime of electronic devices and cause malfunction and performance degradation. In particular, in the case of a display device having a light source that generates a large amount of heat, the accumulated heat is a major cause of deterioration of the display quality of the display device. As electronic devices such as display devices and portable terminals become more sophisticated and miniaturized, electronic parts included in the electronic devices are becoming highly integrated, which, in turn, increases heat density. Therefore, there is a need for a technology that can effectively remove generated heat. As an example method of removing the heat, a heat radiation sheet including a heat radiation member may be placed adjacent to an electronic part that generates heat.
{ "pile_set_name": "USPTO Backgrounds" }
This type of tire for use in dump trucks and the like is used in an environment where the road surface is in poor condition such as a construction site or a mine. Accordingly, contact with pebbles, rocks, etc. on the road often causes a deep cut in the tread portion. In the case where the tread portion receives such a deep cut, the separation between the tread rubber and the belt located on the tire radial inner side of the tread rubber is more likely to develop in the tire circumferential direction from the initial cut when the tire circumferential shearing strain in the tread portion during tire use is greater or the tread temperature is higher. For example, Patent Literature (PTL) 1 discloses the following tire to suppress the heat generation of the tire center portion. A center block row is formed in the tread portion by circumferential grooves and lateral grooves, and each block constituting the center block row has a block groove one end of which is open to the lateral groove adjacent to the block on one side in the tire circumferential direction and the other end of which is open to the lateral groove adjacent to the block on the other side in the tire circumferential direction.
{ "pile_set_name": "USPTO Backgrounds" }
The present invention relates generally to the formation and development of charge patterns and more particularly to the formation of imagewise non-uniform charge patterns and the development thereof with finely divided marking material. In conventional xerography, a photoconductive surface is uniformly charged in the dark with a charge of one polarity. The charged surface is exposed to a pattern of radiation to which it is sensitive, and the charge is dissipated in the radiation-struck areas. An imagewise uniform charge pattern remains in the non-radiation-struck areas. The imagewise uniform charge pattern is normally developed by contacting the surface with a finely divided, colored toner which carries a charge of the opposite polarity. Because opposite polarities attract, the toner particles adhere to the photoconductive surface in the area of the uniform charge pattern. The toner particles are most usually charged to the opposite polarity prior to development by rubbing contact with a carrier material. The carrier material is one which is removed from the toner material in the triboelectric series. The carrier material is usually in the form of particles of a larger size than the toner particles; although the carrier may, in some cases, be a liquid. The toner is usually applied to the surface by cascading or flowing the toner or a toner-carrier combination (generally referred to as developer) across the surface. Other well known toner application methods include magnetic brush development, electrophoretic development and out-of-contact liquid development, such as that described in U.S. Pat. No. 3,084,043 to Gundlach. Normal xerographic development has met with great commercial success. However, there remain areas where improvement is desirable. For example, photoconductive surfaces useful in most commercial xerographic development should be from about 10 to about 60 microns thick. Such thicknesses can be expensive and complicated to manufacture. Any imaging process which enables the use of thinner photoconductive layers would be an improvement. The toner-carrier combination which is well known in normal xerography is somewhat dependent on the ambient relative humidity for successful operation. The humidity is preferably lower. Proper triboelectric charging of the toner is difficult if the humidity is too high. Another difficulty of the toner-carrier combination is that the carrier can become coated with a thin layer of toner material after long periods of use. This is generally referred to as carrier aging. Such coated carrier material cannot be used efficiently to triboelectrically charge the toner material. An imaging process which enables the use of a toner material which does not have to be charged to one polarity or another before development is also desirable. A toner material which is readily useful without a carrier material would also be an improvement. In normal xerography, the toner particles adhere to (develop) the photoconductive surface at the point of charge differential. For example, in normal xerography a plate is charged to about 1,000 v. and then imagewise exposed. Exposure reduces the charge in the light struck areas to about 200 v., leaving about 800 v. in non-light struck areas. The line between a 200 v. area and an 800 v. area on a surface attracts toner particles (see FIG. 1). However, solid area coverage of a large area of uniform 800 v. charge cannot normally be accomplished without the aid of such sophisticated and complex mechanisms as magnetic brush developers or development electrode systems. U.S. Pat. No. 2,777,418 to Gundlach shows a typical development electrode used to achieve said area coverage of a large uniform charge pattern using charged toner. A development system which would make available solid area coverage without such complex mechanisms and with uncharged toner is desirable. Even when magnetic brush development and a developer electrode are used to achieve solid area development, the problem of "developer starvation" is observed. This undesirable phenomenon manifests itself as a reduction of density as large solid areas are developed. The reduction can be quite dramatic and unattractive. It is generally understood to occur because of the limited speed at which the typical toner-carrier type of developer can provide sufficient toner particles of the proper polarity. The development of a charge pattern by a toner particle of one charge leaves a net opposite charge on the carrier. This results in the carrier attracting the remaining toner with an increased attraction, making it more difficult for the remaining toner to leave the carrier and develop subsequent charge patterns. Normally this undesirable situation can be remedied only by replenishing the carrier with toner. A marking method which avoids developer starvation would be useful. In normal xerography, a developed image must oftentimes be transferred to a receiving sheet if it is to be useful. Such transfer is a critical operation which must be handled carefully and with great control to achieve complete transfer while avoiding smearing. There are normal xerographic methods which avoid transfer by coating a photoconductive layer on a conductive paper and developing the image directly on the coated paper. However, these methods require conductive papers and expensive coating treatments during manufacture. They also often lend themselves to liquid (bath) development which can be relatively slow and which sometimes can produce damp copies having an unpleasant odor. An imaging system which can avoid the difficulties associated with the normal xerographic transfer step is desirable. An imaging system which enables development directly onto the final copy while avoiding the need for a photoconductive coating on the final copy and the need for liquid development also would be useful. Other image development systems have been achieved which do not overcome these disadvantages. For example, in U.S. Pat. No. 3,318,698, F. A. Schwertz discloses a means for creating a charge pattern on an insulating surface. The surface is first frosted in an imagewise pattern and then uniformly charged. The frosted areas retain less of a charge than do the unfrosted areas, and a charge differential is created between the charged and uncharged areas. The charge differential of Schwertz is developable by well known xerographic methods. However, because the charge pattern of Schwertz is all of one polarity, it requires a toner which is triboelectrically charged to the opposite polarity. The system disclosed by Schwertz also has the solid area coverage problems discussed above in connection with charge differential development. In U.S. Pat. No. 3,043,217 to L. E. Walkup, there is disclosed a development system in which a charge pattern is formed on an insulating surface and is developed with a typical xerographic developer. Although this system has many uses, it also shares many of the disadvantages of the Schwertz method. In U.S. Pat. No. 3,250,636 to R. A. Wilferth, a magnetic imaging system is disclosed. In that system, a non-uniform pattern of magnetic microfields is established in a magnetizable layer. The uniform pattern is selectively removed by Curie point erasure, leaving an imagewise pattern of magnetic microfields. Curie point erasure is a well known technique and comprises heating a magnetized material above a known critical temperature at which its molecules become disoriented and the material loses its magnetic properties. Curie point erasure is sometimes accomplished by such techniques as flash heating a magnetized material with a Xenon flash lamp while protecting the image area with a mask. While the technique of Wilferth avoids the solid area coverage problems of the prior art, it requires a magnetizable imaging layer and magnetically attractable toner particles and it is not compatible with well known optical imaging methods. Also, it is generally limited to forming dark images because magnetically attractable toners are most usually of a rust or black color. Maksymiak, in U.S. Pat. No. 3,759,222 discloses the use of a non-uniform charge pattern on a transfer member to transport magnetic toner particles to an imaging member which carries a magnetic image. The transfer member of Maksymiak comprises a conductive drum coated with a thin dielectric layer on which is supported a conductive screen. A potential difference is established between the screen and the drum so that a non-uniform charge pattern exists over the surface of the drum. The transport member of Maksymiak does not provide an imagewise non-uniform charge pattern and does not overcome the difficulties of magnetic imaging pointed out above. The existence of field gradients at the edges of xerographic charge patterns and in periodic xerographic charge patterns has been disclosed by H. E. J. Neugebauer (Appl. Opt. 3, 385 (1964) and R. M. Schaffert, Phot. Sci. Eng. 6, 197 (1962). However, the use of conductive toners to attempt to develop such field gradients results in discharge and loss of the image. The use of charged insulating toners or uncharged insulating toners results only in edge development, as described above and in connection with FIG. 1, below.
{ "pile_set_name": "USPTO Backgrounds" }
Typically, electrical equipment, such as, for example, a cordless power tool, is powered by a rechargeable battery. The battery may be periodically charged in a compatible battery charger.
{ "pile_set_name": "USPTO Backgrounds" }
A plasma supply device is a plasma power supply that supplies plasma processes with power. The plasma supply device operates at a basic frequency that, when used as a plasma power supply, should only deviate slightly from a theoretical value. Typical basic frequencies are, for example, 3.39 MHz, 13.56 MHz, 27 MHz, 40 MHz, and 62 MHz. The inverter, which has at least one switching element, generates from the DC signal of a DC power supply an alternating signal that changes its sign periodically at the rate of the basic frequency. For this purpose, a switching element is switched backwards and forwards between a conducting and a non-conducting state within the cycle of the basic frequency. An output network generates from the alternating signal generated by the inverter a sinusoidal output signal at essentially the predetermined basic frequency. A plasma is a special aggregate condition that is produced from a gas. Every gas consists in principle of atoms and/or molecules. In the case of a plasma, the gas is largely ionized, which means that the atoms and/or molecules are split into positive and negative charge carriers, i.e., into ions and electrons, due to the supply of energy. A plasma is suitable for machining workpieces because the electrically charged particles are highly reactive chemically and can also be influenced by electrical fields. The charged particles can be accelerated by means of an electrical field on a workpiece, where they can release individual atoms from the workpiece on collision. The released atoms can be removed by gas flow (etching) or coated on other workpieces (production of thin films). A plasma can be used to machine extremely thin layers, for example, in the region of few atom layers. Typical applications for plasma machining are semiconductor technology (coating, etching, etc.), flat screens (similar to semiconductor technology), solar cells (similar to semiconductor technology), architectural glass coating (heat protection, dazzling protection, etc.), storage media (CD, DVD, hard discs), decorative coatings (coloured glasses, etc.), and tool hardening. These applications impose high demands in terms of accuracy and process stability. To generate a plasma from a gas, energy is supplied to the gas. Energy can be generated in different ways, for example, with light, heat, or electrical energy. If energy is generated using electrical energy, then the plasma is ignited with the electrical energy. A plasma for machining workpieces is typically ignited in a plasma chamber, for which purpose an inert gas, e.g., argon, is generally conducted into the plasma chamber at low pressure. The gas is exposed to an electrical field that is produced by electrodes and/or antennae. A plasma is generated or is ignited when several conditions are met. A small number of free charge carriers must be present, and in most cases, use is made of the free electrons that are always present to a small extent. The free charge carriers are accelerated so much by the electrical field that they release additional electrons when colliding with atoms or molecules of the inert gas, thus producing positively charged ions and even more negatively charged particles (electrons). The additional free charge carriers are again accelerated and on collision produce additional ions and electrons. An avalanche effect is created. The natural recombination counteracts the constant generation of ions and electrons, i.e., electrons are attracted by ions and recombine to form electrically neutral atoms and/or molecules. Therefore energy is constantly supplied to an ignited plasma in order to maintain it. Plasma power supplies are used for generating or igniting and maintaining a plasma, but can also be used for exciting gas lasers. Plasma power supplies have the smaller dimensions to ensure that they can be arranged in the application close to the plasma discharges. They should have the highest possible repeat accuracy and operate precisely, with the smallest possible losses to achieve high efficiency. Further requirements are minimal production costs and high maintenance friendliness. If possible, plasma power supplies are provided without mechanically driven components, and fans can be undesirable because of their limited life and the risk of contamination. Furthermore, plasma power supplies should be as reliable as possible, should not overheat, and should have a long operating time. Due to the high dynamics and often chaotic conditions in plasma processes, a plasma power supply is subject to much more stringent requirements than any other power supply. An un-ignited gas, which has only a very small number of free charge carriers, has an almost infinitely high impedance. Because of its large number of free charge carriers, a plasma has a relatively low impedance. When the plasma is ignited, therefore, there is a rapid impedance change. Another characteristic of an ignited plasma is that the impedance can vary very quickly and often unpredictably, and the impedance is then said to be dynamic. The impedance of the plasma is still non-linear to a great extent, which means that a variation in the voltage on the plasma does not correlate to a similar variation in current. For example, the current can increase much more quickly as the voltage increases due, for example, to an avalanche effect, or the current can also decrease as the voltage increases at so-called negative impedance. If a power supply discharges a power in the load direction, such as a plasma load, which flows at finite speed towards the load, but cannot be absorbed there because the same current is not set when the voltage is present on the load due to the different impedance, only that proportion of the power that is calculated from voltage and current to obtain the load is absorbed, the remaining proportion of the power being reflected. In fact this also takes place in power supplies with low frequencies, and also in direct current, but only in the latter does it take place so quickly that the voltage at the output of the power supply has in practice not yet changed by the time the reflected energy returns. To the observer, therefore, this happens simultaneously. However, in high frequency technology with frequencies above around 1 MHz, the voltage and current at the output of the power supply have generally already changed by the time the reflected power returns. The reflected power has a considerable influence on the power supplies in high frequency technology. Reflected power can destabilize power supplies and prevent the supply systems from operating according to the regulations. Because of incorrect adaptations, the reflected power only has proportions of the basic frequency at constant impedances. The reflected power cannot be blocked or absorbed with filters because filters cannot distinguish between forward (to the load) running waves and backwards (from the load) running waves, and would consequently also block or absorb the forward running waves. In order to reduce or minimize the reflected power, so-called impedance adapter elements or networks are used. Impedance adapter elements or networks can be produced using high frequency technology by combinations of inductances, capacitances, and resistances, with resistances not being absolutely necessary. However, if the load is not a constant impedance, but is a dynamic and non-linear impedance, at least two additional problematic phenomena can arise. First, energies can be generated by the non-linear, dynamic impedance at frequencies that differ from the basic frequency, and proportions of these frequencies are conducted in the direction of the power supply. These are blocked or absorbed by filters. Second, the impedance adapter elements cannot follow the fast dynamic impedance variations sufficiently quickly, thus giving rise increasingly to reflections at the basic frequency, which reflections are conducted by the dynamic impedance to the power supply. Unlike in other power supply systems, plasma power supplies need to be able to be loaded with any incorrect termination, from no load through short-circuit, from infinitely high capacitive load to infinitely high inductive load. At any point on the Smith graph, a plasma power supply must be able to supply power for at least a short period of time and must not suffer permanent damage in doing so. This is linked to the high dynamics and the often chaotic conditions in a plasma process. In addition, frequencies within a wide range and differing from the basic frequency can occur, and these frequencies should be prevented from causing permanent damage to the plasma power supply. The detection and rapid disconnection of an incorrect terminal are allowed in this case, but the plasma power supply should not be damaged if at all possible.
{ "pile_set_name": "USPTO Backgrounds" }
This invention relates to "milled" tooth rotary cone rock bits and methods of manufacture therefor. Rotary cone rock bit s for drilling oil wells and the like commonly have a steel body which is connected to the bottom of a long pipe which extends from the earth's surface down to the bottom of the well. The long pipe is commonly called a drill string. Steel cutter cones are mounted on the body for rotation and engagement with the bottom of the well being drilled to crush, gouge, and scrape rock thereby drilling the well. One important type of rock bit, referred to as a milled tooth bit, has roughly triangular teeth protruding from the surface of the cone for engaging the rock. The teeth are typically covered with a hard facing material harder than steel to increase the life of the cone. The teeth are formed into the steel cone by material-removal processes including turning, boring, and milling. Thus, the cone is referred to as a milled tooth rock bit cone because the teeth are manufactured by milling the teeth into a forged steel preform. The cones may also be referred to as steel tooth cones because they are predominantly manufactured from steel. A milled tooth rock bit cone can have 69 or more milled surfaces, five or more bores, and three or more turned surfaces. Thus, the production of a milled tooth rock bit cone is a labor intensive process, and a majority of the cost of a milled tooth rock bit cone is attributable to the labor cost. The cost is also increased by the waste of raw material which is machined away during the material removal process. The machining processes also leave sharp edges and corners on the finished cone. The sharp edges tend to crack, and the cracks propagate through the cone and through the hard facing, reducing the useful life of the cone. The sharp corners are plagued by stress concentrations which also promote cracking of the cone. Thus, teeth geometry must be limited to avoid sharp edges and corners. Further, the geometry of the teeth is limited by the capability of the milling process making infeasible some tooth shapes that increase the rate of penetration without breakage. To address these limitations, some powder metallurgy techniques have been suggested to manufacture "milled" tooth rock bit cones. For instance, one process currently used utilizes a pattern to form a flexible mold which is filled with powdered metal. The mold is cold isostatically pressed to partially densify the powdered metal. Isostatic pressure is pressure equally applied on all sides of the mold. The partially densified part, called a green part or preform, is then heated and rapidly compressed to full density by a quasi-isostatic process. To create the preform, the powdered metal, usually steel, is poured into the flexible mold while the mold is vibrated. Vibrating the mold during filling uniformly packs the powder in the flexible mold. The flexible mold is supported during the cold isostatic pressing by tooling which allows the deformation necessary to compress the pattern. After the mold is compressed, the preform is removed from the mold and subjected to uniform heating. Once the preform is heated, it is transferred to a central position in a cylindrical compression cavity in which it is surrounded by a bed of granular pressure transfer medium heated to approximately the same temperature as the preform. The pressure transfer medium is then axially compressed creating a quasi-isostatic pressure field acting on all surfaces of the preform. The radial pressure acting on the preform approaches a theoretical maximum of one-half of the axial pressure acting on the preform. After compression, the part is removed from the cavity and allowed to cool slowly over a two (2) hour time period. This powder metallurgy process requires two compression steps, and because the non-isostatic compression step causes a non-uniform reduction in size of the preform, its pattern is complex. The second compression process is essentially a hot pressing process, which is expensive and inefficient but only one part can be made at a time. Further, the steps required to prepare the part for the hot pressing process are complex and time consuming. Then, the process is not economical. Other powder metallurgy process including powder injection molding have been utilized to fabricate small parts. In summary, this process begins by pelletizing or granulating a mix of powder metal and binder before injecting the pellets or granules into the mold. The mold is then removed, and the part is debinded and sintered. This process has only been utilized for small parts with thin cross-sections and heretofore has not been utilized for the production of milled tooth rock bit cones. Thus, reduction in the required labor to fabricate a "milled" tooth rock bit cone is desirable to enhance the production rate and reduce production cost of the milled tooth rock bit cone. It is also desirable to diversify the geometric shapes of the teeth to increase the rate of penetration without the need for complexly shaped molds and preforms. Thus, the successful application of powder injection molding to produce "milled tooth rock bit cones" is desirable to bring about such an increase in the rate of penetration and decrease in the cost of rock bits which translates directly into reduction of drilling expense.
{ "pile_set_name": "USPTO Backgrounds" }
For a common slide door opening and closing device, an upper rail, a lower rail and a center rail are attached to an upper part, a lower part and a rear-center part of an entrance opening at a vehicle side face, respectively. To make a slide door opening and closing operation, rollers, rotatabily supported to brackets fixed to the slide door, are slid along each rail. In an automatic opening and closing device for a vehicle slide door, a belt driving system, a wire driving system, a cable driving system or others is commonly applied with being assembled on vehicle step panel which is used as an entrance step for getting on or off of passenger. An example of a belt driving system has been disclosed in a Japanese Patent Laid-Open Publication published as No. Hei. 10(1998)-8828. In the disclosed belt driving system, a motor, a driving pulley driven by the motor, a driven pulley (provided at both fore side and rear side of a step panel), and an idle pulley are displaced on the under surface of the step panel and assembled as a belt driving system module. The step panel is attached on the vehicle body with a space under it for accommodating the belt driving system mechanism therebetween. A loop-like timing belt extends fore and aft and is wound on each pulley, and a part of a slide door is fixed to the timing belt through a bracket. Then, the slide door can move to open or close. In this system, an output rotation axis of a drive motor points to vertical direction because of the relationship between the driving pulley and the timing belt. In this configuration, projected area on a horizontal plane of the drive motor is wider than projected area on a vertical plane, because the drive motor usually is attached with drive mechanisms including such as a magnetic clutch. In actual design, a part of the step panel has to be deformed to inflate upward to ensure the height space for accommodating the drive motor. Therefore, the wider the infrared area of the step panel for the space to the drive motor becomes, the narrower the entrance step area for the passenger is. Therefore, design flexibility of the step panel is restricted, and ease in getting on and off becomes worse. For a wire or a cable driving system, Japanese Patent Laid-Open Publications published as Nos. Hei. 8(1996)-232539 and Hei. 5(1993)-61432 have been disclosed. In these systems, a driving pulley driven by a motor and a driven pulley (provided at both fore side and rear side of a step panel) are provided under the step panel of a vehicle. In addition, a loop-like wire (or cable) is wound on each pulley, and a part of a slide door is fixed to a one end of the wire. Then, the slide door can move to open or close. These wire or a cable driving system have design flexibility for the arrangement of an output rotation axis of a drive motor points in vertical or horizontal direction. However, an effective length of the wire is considerably changed, when a lower roller of the slide door is moving at either a straight line path or a curved line path of a rail. Then, a tension adjusting pulley is required to provide on the path of the wire. It makes the structure complex and makes the number of parts increase. In the structure of a cable driving system disclosed in Japanese Patent Laid-Open Publication published as No. Hei. 8(1996)-232539, durability of the cable may be worse, because sliding operation between the cable and a guiding tube may cause wearing. These difficulties can not be eliminated even if the belt is simply alternated to a wire or a cable for the belt driving system described previously. More specifically, even if the belt is simply alternated to a wire, the conversion of the output rotation axis direction of the driving motor and the change of the effective length of the wire can not be eliminated all together.
{ "pile_set_name": "USPTO Backgrounds" }
Linear programming is a technique of determining an optimal solution to a set of defined constraints, which are represented as linear equalities or inequalities. Different classes of linear programs exist, including integer linear programs (IPs or ILPs) and mixed integer linear programs (MIPs or MILPs). Integer linear programs are those in which decision variables have to take integral values and satisfy the set of constraints, while MILPs have some decision variables that need to take integral values and some decision variables that can take on fractional values.
{ "pile_set_name": "USPTO Backgrounds" }
The present invention is related to the area of chemical synthesis. More specifically, one embodiment of the present invention provides certain photolabile compounds and methods for their use as photocleavable linkers. The use of a photolabile molecule as a linker to couple peptides to solid supports and to facilitate the subsequent cleavage reaction has received considerable attention during the last two decades. Photolysis offers a mild method of cleavage which complements traditional acidic or basic cleavage techniques. See, e.g., Lloyd-Williams et al. (1993) Tetrahedron 49: 11065-11133. The rapidly growing field of combinatorial organic synthesis (see, e.g., Gallop et al. (1994) J. Med. Chem. 37: 1233-1251; and Gordon et al. (1994) J. Med. Chem. 37: 1385-1401) involving libraries of peptides and small molecules has markedly renewed interest in the use of photolabile linkers for the release of both ligands and tagging molecules. A phenacyl based linking group (see 1 below) has been described. See Wang, (1976) J. Org. Chem. 41: 3258. ##STR1## An ortho-nitrobenzyl support (see 2 below) derived from 4-bromomethyl-3-nitrobenzoic acid has been widely employed as a photolabile support for the generation of both peptide acids and amides. See Rich et al. (1975) J. Am. Chem. Soc. 97: 1575-1579 and Hammer et al. (1990) Int. J. Peptide Protein Res. 36: 31-45. ##STR2## Photolabile support 2, though useful, does suffer from several limitations. For example, workers have been unable to obtain high yields of methionine-containing peptides from the support without substantial contamination with methionine sulfoxide. See Rich supra and Hammer supra. One solution has been to employ methionine sulfoxide throughout the peptide assembly and to subsequently reduce back to methionine to avoid any ambiguities associated with partial oxidation (see, Lloyd-Williams et al. (1991) J. Peptide Protein Res. 37: 58-60 and Lloyd-Williams et al. (1993) Tetrahedron 49: 10069-10078), but this clearly detracts from the usefulness of the technique. This support also suffers from unduly slow cleavage kinetics, with typical photolysis times in organic solvents ranging from 12 to 24 hours. Moreover, photolysis of the support generates a reactive and chromogenic nitroso-aldehyde on the support which can trap liberated compounds and may act as an internal light filter to slow the rate of cleavage. See Patchnornik et al. (1970) J. Am. Chem. Soc. 92: 6333-6335. Pillai and co-workers have described an .alpha.-methyl-ortho-nitrobenzyl support designed to eliminate formation of the nitroso-aldehyde, but they observed inefficient release of peptides longer than five residues due to poor swelling of the resin. See Ajayaghosh et al. (1988) Tetrahedron 44: 6661-6666. In the course of optimizing the photolithographic synthesis of both peptides (see Fodor et al. (1991) Science 251: 767-773) and oligonucleotides (see Pease et al. (1994) Proc. Natl. Acad. Sci. USA 91: 5022-5026, we had occasion to explore the use of a variety of ortho-benzyl compounds as photolabile protecting groups. See PCT patent publication Nos. WO 90/15070, WO 92/10092, and WO 94/10128; see also U.S. patent application Ser. No. 07/971,181, filed Nov. 2, 1992, and Ser. No. 08/310,510, filed Sep. 22, 1994; Holmes et al. (1994) in Peptides: Chemistry, Structure and Biology (Proceedings of the 13th American Peptide Symposium): Hodges et al. Eds.; ESCOM: Leiden; pp. 110-12, each of these references is incorporated herein by reference for all purposes. Examples of these compounds included the 6-nitroveratryl derived protecting groups, which incorporate two additional alkoxy groups onto the benzene ring. Introduction of an .alpha.-methyl onto the benzylic carbon facilitated the photolytic cleavage with >350 nm UV light and resulted in the formation of a nitroso-ketone. Photolabile amide protecting groups for C-termini of peptides which rely on the same basic ortho-nitro benzyl chemistry have also been reported. See Henricksen et al. (1993) Int. J. Peptide Protein Res. 41: 169-180; Ramesh et al. (1993) J. Org. Chem. 58: 4599-5605; Pillai (1980) Synthesis 1-26; and Pillai et al. (1979) Tetrahedron Lett. 3409-3412. See also Bellof and Mutter (1985) Chimia 39: 10. A photocleavable linker should be stable to variety of reagents (e.g., piperidine, TFA, and the like); be rapidly cleaved under mild conditions; and not generate highly reactive byproducts. The present invention provides such linkers.
{ "pile_set_name": "USPTO Backgrounds" }
The popularity of online social communities has grown at an incredible rate. The fast growing popularity reflects a large population of active users. A large population of active users can be used to distribute information efficiently and quickly. Security concerns prevent or substantially limit access to online social communities, thus preventing or substantially limiting access to the utility and entertainment of social network websites.
{ "pile_set_name": "USPTO Backgrounds" }
The present invention relates generally to devices and methods for detection of analytes in test samples. More specifically, the present invention provides solid phase test devices and methods that combine an internal indicator on the test with an external mark located on a support. Various analytical procedures and devices are commonly employed in detection assays to determine the presence and/or amount of substance of interest or clinical significance which may be present in biological or non-biological fluids. Such substances are generally termed “analytes” and can include antibodies, antigens, drugs, or hormones. The present invention includes, but is not limited to, lateral flow chromatography assay formats. Generally, these assays have an extended base layer on which a differentiation can be made between a sample application region and an evaluation region. In typical use, the sample is applied to the sample application region, flows along a liquid transport path which runs parallel to the base layer, and then flows into the evaluation region. A capture reagent is present in the evaluation region, and the captured analyte can be detected by a variety of protocols to detect visible moieties associated with the captured analyte. For example, the assay may produce a visual signal, such as color change, fluorescence, luminescence, and the like, when indicating the presence or absence of an analyte in a biological sample. Vertical flow devices and assays are also contemplated herein. Generally, these assays, similar to those described in U.S. Pat. No. 4,632,901, involve the introduction of a liquid sample to a device and allowing the fluid sample to pass through one or more layers to generate a result. Results, similar to the assays described above and below, may be in the form a visual signal. Optimally, such test devices will provide a characteristic signal when the analyte is present in a sample, and a different signal when the analyte is absent from a sample. Most typically, the test device will display a “plus” (+) signal in the presence of analyte, and a “minus” (−) signal in the absence of analyte. The plus/minus test result format has enjoyed enthusiastic customer response and wide commercial success. Test devices of this kind are well known in the art, and play an important role in areas such as clinical chemistry. They are used by skilled clinicians and lay person alike. Thus, there is a strong impetus to provide devices that are simple and reliable. Desirably, the assays are single-step devices wherein the user need only apply the sample prior to viewing the result. Single-step devices obviate the necessity of performing complicated and time consuming processing steps that may introduce errors in the end result. Examples of such assays include pregnancy tests, ovulation tests, various urine, saliva, spinal, and blood tests, as well as other clinical or diagnostic assays. Test devices typically use materials that specifically bind with an analyte of interest. A homologous pair of specific binding pair members (“sbp members”), usually an immunological pair comprising a ligand and a receptor (antiligand), is involved, wherein one of the sbp members is labeled with a label that provides a detectable signal. The immunoassay methodology results in a distribution of the signal label between signal label bound in a complex of the sbp members and unbound signal label. The differentiation between bound and unbound signal label can be a result of physical separation of bound from unbound signal label or modulation of the detectable signal between bound and unbound signal label. In developing an assay device, there are many considerations. One consideration is to provide substantial differentiation between the observed signal resulting from signal label when bound as compared to unbound. Another consideration is the ease with which the observed signal can be detected and serve to differentiate between the presence or absence of analyte of interest. Other factors include the precision with which the test devices must be manufactured. In factoring this consideration it is important to include registration or indexing capabilities in vertical flow test devices of the present invention. These capabilities, as described in detail below, are important for testing accuracy, reproducibility and ease of use and reading results. Therefore, in developing an assay that can be used by untrained personnel, such as assays to be performed in the home, medical offices and the like, the technique for performing the assay should be simple, and the method of manufacturing the assay should be straightforward, Plus/Minus Assays Of particular interest to the present invention are test devices of the type described in U.S. Pat. No. 5,145,789 to Corti et al., the disclosure of which is incorporated herein by reference. Corti et al. discuss a built-in positive control to indicate successful operation of a pregnancy test device. The positive control is envisaged as a horizontal tract that always stains, independent of the presence of hCG in the urine, and is described as an area on a membrane that contains immobilized hCG. Regardless of whether hCG is present in the biological sample, it is intended that during operation, the upstream mobile labeled hCG binding reagents will always bind to the immobilized hCG, thereby forming a horizontal line, or minus sign, in the reading area. A similar approach for providing a minus sign in a test device is described in U.S. Pat. Nos. 4,916,056, 5,008,080 and 5,160,701 to Brown, III et al., the disclosures of which are incorporated herein by reference. As illustrated, the positive control is formed by providing a binding substance within the test strip matrix, and is formed in the shape of a rectangular bar, or minus sign. The binding substance of the minus sign is intended to bind the labeled material regardless of the presence or absence of the analyte of interest in the test sample. Another approach for providing a positive control in a test device is described in EP Patent Publication No. 0 249 418 to Graham, Jr., the disclosure of which is incorporated herein by reference. As described, the control zone has anti-human IgG or IgM immobilized thereon, for nonspecifically capturing human immunoglobulin ubiquitously present in all similar human aqueous samples. The immobilized antibody is intended to provide a signal in a “minus” pattern, regardless of the presence or absence of the analyte of interest in the test sample. Osikowicz et al., in U.S. Pat. No. 5,075,078, describe yet another approach for providing a positive control in a plus/minus test device. The positive control is disposed on a test strip in a rectangular bar configuration. The control bar is oriented on the strip so that it lies neither perpendicular nor parallel to the direction of fluid flow, but rather lies at an intermediate orientation, i.e., at a 45 degree angle. Still yet another approach for providing a positive control in a test device is provided in U.S. Pat. No. 5,401,667 to Koike. As described, the test device provides a plus/minus format, but considers alternative geometric symbols as well. A portion of the chromatographic medium is removed, or otherwise partially blocked, thereby affecting the flow-path of the liquid. It is suggested that this modification enhances the signal of the device. Wong et al., in EP Patent No. 0 260 965, describe another test device that utilizes the plus/minus format. Wong et al. discuss a multiple-step diagnostic assay with a horizontal positive control line sprayed onto a test membrane. The previous methods discussed above accomplish the “appearance” of a minus sign (−) by placing an indicator (positive control) line perpendicular to the test line, directly onto the test strip. Typically, the control line develops with any sample flow, while the test line develops only with a positive sample flow. Thus, the previous assays involve a control mechanism inherent to the matrix membrane test strip, and require a specific manufacturing step to apply the control line to the strip. Other previous devices display a printed minus sign positioned on the matrix and across the test line. These devices typically incorporate a positive control line downstream from, and parallel to the test line. Such devices are limited as the test strips may present a line that is visible before the sample is added. Previous methods are further disadvantaged as the additional manufacturing step involves a difficult placement procedure to orient the perpendicular line directly in the center of the viewing window. Whether the perpendicular line is a printed minus sign, or a reagent-based control line, this approach is particularly ill suited for certain matrix construction procedures, including web processing methods that involve a continuous flow or continuous roll application approach. Therefore, it would be desirable to provide a test device that does not require this extra processing step of depositing a perpendicular line onto the test strip, or does not leave a line that is visible before the sample is added to the device. This invention fulfills these and other needs. Transparent Membranes The use of transparent test strips in diagnostic assays is known in the art. In U.S. Pat. No. 4,824,640, Hildebrand et al. discuss a transparent reagent carrier layer suitable for evaluation by transmission photometry. As described, the transparent nature of the film of plastic provides a suitable carrier material as compared to opaque films. The use of a transparent test strip is also discussed in U.S. Pat. No. 5,110,550 to Schlipfenbacher et al. As described, this test device includes a white non-transparent covering layer situated above a color-forming layer. During operation of the test device, the covering layer becomes transparent in the moist state. Through the transparent covering layer, the user is then able to observe any reaction occurring in the color-forming layer below. The use of a clearing agent in an immunochromatographic assay is discussed in U.S. Pat. No. 6,165,798 to Brooks. As described, the test strip membrane is rendered transparent by wetting the membrane with a clearing agent, thus reducing the amount of light scattered by the membrane fibers. In U.S. Pat. No. 6,187,268. Albarella et al. describe a transparent flow through membrane for use in test devices, but do not suggest a control feature to indicate a positive or negative test result. The membrane described in Albarelia et al. is not configured to become transparent only when wet. The membrane is transparent whether wet or dry. While conceivably workable in some circumstances, the previous detection systems that employ transparent membranes are of limited utility. There is no teaching or suggestion in current art for a test device with a transparent membrane that utilizes a control feature to indicate a positive or negative test result as provided by a mark on the underlying support. In view of the foregoing, there remains a need in the art for a simple, efficient method for adding a positive control to a solid phase assay that does not require the manufacturing step of fixing a positive control binding member to the assay test strip, and that does not leave a substantially visible signal before the sample is added to the device. It would further be desirable to achieve improved test device formats that incorporate transparent membranes as part of a control or display feature. Additionally, the assay of the present invention should overcome the disadvantages described above in connection with the previous test device systems.
{ "pile_set_name": "USPTO Backgrounds" }
Over the past few years, dramatic increases in the cost of materials and labor have focused attention upon methods of reducing costs in various types of industrial and commercial devices. In the area of power transmission, it is usual that complex structures will be utilized to satisfy a relatively low torque transmission requirement. Recently industry has been concerned with production of transmission assemblies including utilization of stamped components. These structures, however, continue to be of complex construction including numerous components. Consequently, any malfunction within the device necessarily requires substantial time and possible expenditure of substantial sums toward replacement parts in putting the device back into an operable condition. Change speed mechanisms currently used in bicycles are a specific area where the speed-changing device is of a complex nature. The usual three-speed, ten-speed, twelve speed, or eighteen speed shfting structure includes a plurality of sprockets in combination with a relatively sensitive shifting mechanism. The number of individual components in these devices is likewise substantial and repair necessarily involves expenditure of considerable time and money toward replacing a deflective component somewhere in the mechanism. A further disadvantage resides in the fact that these bicycle change speed devices can only be operated while the bicycle is in motion, thus preventing a change in a stopped or static condition.
{ "pile_set_name": "USPTO Backgrounds" }
Lasers are used in many different industries for many different purposes, such as, for example, in the medical industry for medical procedures, in the printing industry in laser printers, in the defense industry in a variety of defense applications, and in the optical communications industry for transmitting and receiving optical signals. In many applications, the output power of the laser is monitored and controlled to maintain the output power at a desired or required level. Due to the wide variations in laser parameters such as, for example, laser slope efficiency (SE) and laser threshold current (ITH), maintaining the optical power at a particular level is challenging because temperature and process variations and aging of system components also cause the output power level to vary. Many techniques and systems have been used or proposed to control and maintain the output power of the laser at required levels over temperature and process variations and time. It is common practice in the optical communications industry to use a monitor photodiode to detect light output from a rear portion of the transmitter laser (or a portion of the output power reflected back through optical lenses) and to use this optical feedback to measure and control the average transmitted output power level of laser. In general, the average transmitted output power level, PAVG, of the laser can be controlled by controlling the bias current, IBIAS, of the laser. Thus, if the optical feedback indicates that PAVG has fallen below the required level, increasing IBIAS by an appropriate amount will raise PAVG to the required level. Similarly, if the optical feedback indicates that PAVG has risen above the required level, decreasing IBIAS by an appropriate amount will lower PAVG to the required level. As the optical feedback path described above is used to maintain PAVG at the required level, the laser is modulated with a modulation current, IMOD, to cause the laser output power level to be adjusted between a level, P1, that represents a logic 1, and a power level, P0, that represents a logic 0. The amplitude of the modulation current IMOD1 corresponds to an output power level of P1. The amplitude of the modulation current IMOD0 corresponds to an output power level of P0. The laser threshold current ITH has an amplitude that is sufficient to cause the laser to begin producing laser action (i.e., to emit stimulated radiation). The amplitude of the threshold current ITH needed to produce laser action varies due to factors such as, for example, temperature and aging. Due to these variations in the amplitude of ITH that is needed to produce lasing and the slope efficiency of the laser, adjustments must be made to the amplitudes corresponding to IMOD1 and IMOD0 in order to maintain P0 and P1 at the necessary respective output power levels. A variety of techniques have been used to control the amplitude of the modulation current. One known technique sets the amplitude of the modulation current at a level that achieves a desired extinction ratio (ER) or optical modulation amplitude (OMA) at a fixed temperature or time. The amplitude of the modulation current is then increased or decreased based on an analog temperature coefficient, or in a digital control system, based on a temperature measurement and/or and aging timer. This technique generally provides suitable results if the laser SE variation is controlled well enough to maintain the ER/OMA and the laser performance within specifications. A disadvantage of this technique is that using a single temperature reference point means the adjustment to the amplitude of the modulation current is essentially based on a “guess” of changes of the SE in direction and amount based on statistical or other data. Because the change in the SE often is not linear and can change from positive to negative slope from one temperature to the next, it is difficult or impossible to determine the optimal adjustment in the amplitude of the modulation current. In addition, this technique can also limit laser yields because the SE and ITH limits need to be within sufficiently tight tolerances to guarantee that a suitable level of performance will be achieved without the necessity of testing each laser over temperature and customizing each laser based on the results of testing. Another known technique involves measuring the amplitudes of the modulation current needed to maintain the required output power levels PAVG, P0 and P1 over a range of temperatures on a part-by-part or wafer-by-wafer basis. The amplitude values obtained during testing are programmed into a lookup table (LUT) memory element or other non-volatile memory element. A controller uses a temperature measurement value to index into the memory element and read out the corresponding amplitude value for the modulation current. The amplitude of the modulation current of the laser is then set to the value read out of the memory element. One disadvantage of this technique is that it requires over-temperature testing during manufacturing, which is expensive. In addition, it is difficult to factor in aging when using this technique, which means that the amplitude of the modulation current often will not be set to an optimal level. Another technique that has been proposed involves using a high-speed monitor diode in combination with an amplitude detector to monitor the output power level of the laser and adjust the amplitude of the modulation current to achieve the required output power level. This technique requires an additional high-bandwidth feedback path for the amplitude detector, which increases the cost and complexity of the transmitter. In addition, the performance of the amplitude detector can be significantly affected by temperature variations, which can lead to less than optimal performance. Also, the amplitude detector dissipates a large amount of power relative to the rest of the transmitter. Because of these difficulties, this technique has been proposed, but not actually implemented. Another known technique involves measuring ITH in situ by adjusting the modulation current amplitude while measuring the optical feedback signal to obtain the slope and calculate the corresponding SE. The measured ITH and SE are then used to calculate the amount of modulation current needed. An advantage of this technique is that the existing feedback path for the monitor photodiode is used. A disadvantage of this technique is that it requires a large amount of signal processing to be performed to make the necessary calculations, and is therefore computationally intensive. Another disadvantage is that the method is performed during module power up or module programming, and generally cannot be used while transmitting actual data. A similar technique also uses the existing feedback path to calculate SE, but also modulates a very small amplitude signal at low frequency on top of the IBIAS supplied to the laser. This signal is then extracted from the feedback signal, amplified and used to calculate SE. The calculated SE is then used to determine how to adjust the modulation current amplitude. This method can be used while transmitting actual data, but requires that high accuracy circuits such as amplifiers, analog-to-digital converters (ADCs) and digital-to-analog converters (DACs) be included in the feedback path. In addition, this technique also requires a significant amount of signal processing, and is therefore computationally intensive. A need exists for a way to determine the adjustment needed, if any, to the amplitude of the modulation current in order to achieve a desired or required output power level, which can be used when transmitting actual data, which uses the existing feedback path, which does not require complex high-accuracy circuitry, and which does not require that over-temperature testing be performed.
{ "pile_set_name": "USPTO Backgrounds" }
Wireless mesh networks are gaining popularity because wireless infrastructures are typically easier and less expensive to deploy than wired networks. The wireless mesh networks typically include wired gateways that are wirelessly connected to wireless nodes, or wirelessly connected directly to client devices. Many wireless nodes can collectively form a wireless mesh, in which client devices can associate with any of the wireless nodes. Wireless networks, however, can be more difficult to deploy and maintain than wired networks. That is, wireless networks are typically subjected to environmental influences that make operation of the networks more problematic than wired networks. For example, the wireless links of wireless networks can suffer from fading or multi-path, which degrade the quality of transmission signals traveling through the wireless links. Additionally, wireless networks that include multiple access points can suffer from self-interference. During deployment of wireless networks, it is useful to have information regarding fading, multi-path and self-interference to allow for better deployment. Additionally, this information is useful after deployment to allow determination of the quality of possible network connections, and to provide information that can be used to improve the wireless networks. It is desirable to have a method and apparatus for determining coverage a wireless network, allowing a system operator to make strategic upgrades to the network. It is also desirable to have measured network information available so that potential users can determine a level of expected performance.
{ "pile_set_name": "USPTO Backgrounds" }
The present invention relates to a terminal locking structure in which a terminal lance is prevented from being deformed when a terminal of an electric connector is removed from a connector housing. Heretofore, various structures disclosed, for example, in Japanese Patent Unexamined Publication No. Hei. 3-147282, are known as terminal locking structures of electric connectors. As shown in FIG. 5, in a terminal locking structure 31 of a conventional electric connector 40, a first free end 35a of a flexible terminal lance 35 made of a metal plate is engaged with an engagement portion 33 in a connector housing 32 to thereby lock a terminal 34 at a predetermined position in the connector housing 32. The terminal lance 35 of the terminal 34 is made of a metal plate integrated with a box ceiling plate 36 so as to be bent at a desired angle obliquely backward and outward. A flexible contact plate 39 which cooperates with the box ceiling plate 36 to press and hold a mating contact 50 for the electric connector 40 to thereby electrically connect the maiting contact 50 to the electric connector 40, is bent integrally with a bottom plate 38. An end plate at the first free end 35a of the terminal lance 35 is punched in its center, bent downward and further bent so as to be parallel with the bottom plate 38, so that a second free end 35c is formed. The second free end 35c abuts on a stopper plate 37 supported by a stabilizer 41 of the terminal 34, so that the terminal lance 35 is prevented from projecting out due to its flexibility over a necessary extent. An electric wire caulking portion 42 is provided in the rear portion of the terminal 34 in order to fix an electric wire W and to make electrical connect with a core of the electric wire W. An unlock hole 45 for removing the terminal 34 is provided under a ceiling 44 of the connector housing 32 at the upper position of the box ceiling plate 36 of the terminal lance 35. Further, an insertion hole 47 for leading the mating contact 50 therein is provided in the front wall 46 of the connector housing 32 whereas a slot 49 for preventing the mis-insertion of the terminal 34 is provided in the center portion of the bottom wall 48. In the thus configured terminal locking structure 31, when the terminal 34 with the electric wire W caulked by the electric wire caulking portion 42 in the rear is inserted into the connector housing 32 while being slid along the bottom wall 48 and the slot 49 from the rear of the connector housing 32, the terminal lance 35 projecting upward abuts on a rear end portion 44a of the ceiling 44 so as to be moved down while being pressed. When the terminal 34 is further inserted, the first free end 35a of the terminal lance 35 passes through the engagement portion 33 so that not only the terminal lance 35 is restored to its original state but also the first free end 35a is locked in the front end surface of the engagement portion 33. Accordingly, not only the terminal 34 is locked at a predetermined position in the connector housing 32 but also the front portion of the terminal is positioned by the front wall 46. When the mating contact 50 is inserted into the insertion hole 47 in the front portion of the connector housing 32, the mating contact 50 is held between the contact plate 39 and the box ceiling plate 36 so that the electric wire W is electrically connected to the mating contact 50. In the above conventional terminal locking structure 31, to take out the terminal 34 from the connector housing 32, when an unlock jig J is inserted through the unlock hole 45 so as to be brought into contact with the inclined surface of the terminal lance 35 or the vicinity of the first free end 35a in the direction of the arrow B to thereby press down the terminal lance 35, the first free end 35a is disengaged from the engagement portion 33. Accordingly, when the rear end portion of the terminal 34 is pulled back, the terminal 34 can be removed from the connector housing 32. In this occasion, however, if the inside of the unlock hole 45 cannot be seen so that the unlock jig J is entangled, a root portion 35b of the terminal lance 35 may be deformed by the unlock jig J. Accordingly, there arises a problem that the terminal lance 35 is plastically deformed to spoil its function.
{ "pile_set_name": "USPTO Backgrounds" }
The present invention relates generally to printing techniques and in particular to verifiable printing. Secured paper is commonly associated with currency, stocks, and other financial instruments. The production of conventional secured paper is costly in terms of the paper stock (because specialized paper having security features is required) and the printing equipment (because special inks and printing machinery is required). Many businesses could benefit if secured paper document capability was cost-effectively available. Various techniques are known for uniquely signaturizing a paper document based on inherent characteristics of the paper. This allows a recipient to verify the originality of a document by obtaining a signature of the paper document in question and comparing it against a data store of signatures of documents.
{ "pile_set_name": "USPTO Backgrounds" }
The frequent occurrence of natural and environmental catastrophes (earthquakes, floods, cyclones, aridity/drought, fire catastrophes, etc), war and terror events and associated instability of the financial markets (stock market crashes, etc) have imparted previously unknown importance to risk management and corresponding measures for handling such risk events and catastrophes in recent years for general economic activity, since a high proportion of business volume and a considerable percentage of jobs can be endangered thereby. Particularly in the insurance/reinsurance sector, there has been a long-known backlog in technical automation in many areas. The appearance of the World Wide Web and the resulting possibility of being able to access enormous, decentralized quantities of data have created completely new requirements for industry. The survival of an entire branch of industry may depend, for example, on being able to analyze the relevant data rapidly and reliably in order to be able to take the appropriate measures. This is no longer possible based on human work alone but requires a great deal of automation. This type of automation has an extremely important role for industry and society, for example even in areas with traditionally less technical character, such as the insurance industry, at the present time (when at the same time technical progress is associated with an unavoidable increase in risks and a change in the liability concepts within the definition of third party liability). In Europe, the introduction of the Euro has furthermore inevitably brought new developments which result in greater transparency, facilitate cross-border comparison and the conclusion of cross-border contracts. Because of this, too new possibilities for comparison and automation have resulted which is completely utilized from the prior art up to today only in a few areas. In industry, for expedient risk management for surviving risks, it is essential to know or to be able to estimate, as a boundary condition parameter, a reacquisition value (monetary replacement value) in a country-specific manner in order to be able to determine the necessary capital collateralization levels. With the means known in the prior art, values can only be obtained with considerable uncertainty for the overall risk sum (for example for determining the total insurance sum). This is true not only for industrial plants in so-called developing or emerging countries but also for plants in industrially highly developed countries, such as, for example, Germany, Austria, Italy, Switzerland etc. One of the possible problems is that the company or the object has only an insufficient capital collateralization level for surviving the risk event or is underinsured so that reprocurement is not possible in the event of loss. A further problem with the prior art is that it provides no system or method covering a branch of industry in order to be able to determine such collateralization levels. In other words, coarse estimates or “rules of thumb” are generally used, which can scarcely be automated and with which it is difficult to determine the errors and uncertainties in the calculations. Thus, a considerable uncertainty or inaccuracy with regard to errors always remains in the prior art. A further disadvantage is that technical automation is scarcely possible in the prior art owing to the nonuniform and/or complicated methods. As a result of this, the systems of the prior art do not permit effective dynamic monitoring. The outlay in employee time, costs and material are therefore correspondingly high.
{ "pile_set_name": "USPTO Backgrounds" }
Recently, electronic devices such as servers and personal computers have been remarkably developed in terms of advancements in speed, performance, and the like, and accordingly semiconductor elements such as CPU (Central Processing Unit) used in the electronic devices have been progressively increased in size. As a mounting technology for semiconductor elements, flip chip mounting is known in which a semiconductor element in the form of bare chip is directly mounted on a wiring board with a solder bump. Additionally, to scale up the fine electrode arrangement of semiconductor elements to the electrode arrangement of a wiring board, there is also a mounting method in a BGA (Ball Grid Array) approach in which a semiconductor package having a semiconductor element placed on an interposer is fabricated and mounted on a wiring board with a solder bump interposed therebetween. The semiconductor package for BGA approach is also called a BGA semiconductor package. FIGS. 1A and 1B are cross-sectional views of a BGA semiconductor package 5 in the course of the mounting thereof on a wiring board 1. As illustrated in FIG. 1A, the wiring board 1 has first electrode pads 2 on one main surface thereof. A solder paste 4 is printed in advance on the first electrode pads 2 by screen printing. On the other hand, the semiconductor package 5 includes second electrode pads 6 on a main surface thereof at positions facing the first electrode pads 2. Further, solder bumps 7 are bonded to the upper surfaces of the second electrode pads 6. Then, while the solder bumps 7 are in contact with the solder paste 4, these are reflowed by heating. Thereby, the semiconductor package 5 is mounted on the wiring board 1 as illustrated in FIG. 1B. The surface tension of the solder and the own weight of the semiconductor package 5 determine the shape of the solder bumps 7 after the reflowing, which is normally a drum-like shape bulging at the center as illustrated. Meanwhile, the semiconductor package 5 and wiring board 1 have different thermal expansion coefficients because of the difference in materials. Accordingly, as the semiconductor package 5 generates heat, stress is applied on the solder bumps 7 due to the difference in thermal expansion coefficient. The stress concentrates on portions of the solder bumps 7 where the diameter is the smallest, in other words, around bonded portions A between the electrode pads 2, 6 and the solder bumps 7. As the power supply of the semiconductor package 5 is turned on and off repeatedly, the stress is repeatedly applied to the solder bumps 7 in the bonded portions A. Thus, metal fatigue gradually progresses at the solder bumps 7. Eventually, a crack is generated in the solder bumps 7, and the bonded portions A may be fractured. Patent Literature 1: Japanese Laid-open Patent Publication No. 05-114627 Patent Literature 2: International Publication Pamphlet No. WO 08/114434 Patent Literature 3: Japanese Laid-open Patent Publication No. 2001-118876 Patent Literature 4: Japanese Laid-open Patent Publication No. 08-236898 Patent Literature 5: Japanese National Publication of International Patent Application No. 2005-510618 Non-patent Literature 1: Morita, Hayashi, Nakanishi, and Yoneda, “High Acceleration Test of Lead-free Solder”, 23rd Spring Lecture Meeting of Japan Institute of Electronics Packaging
{ "pile_set_name": "USPTO Backgrounds" }
Not applicable. Not applicable. This invention relates in general to envelope openers and, more particularly, to a slitter-type envelope opener having its handle portion and slitting portion formed as separate components and to a process for making the same. Slitter-type envelope openers open envelopes with considerable ease, are quite compact and are produced inexpensively. They find widespread use in offices throughout the country. The typical opener of this type has a handle which enables one to grip the device and a finger or spear that is located below the handle. Between the handle and spear lies a slot, and at the end of the slot, a blade. The spear is small enough to fit behind the fold of a flap for an envelope and the slot is large enough to accommodate the fold. Thus one, while gripping the opener along its handle, manipulates the spear beneath the flap and then advances it behind the fold. After short distance the blade encounters the fold and slits the envelope along the fold as the device is advanced. Many businesses give the openers away as promotional items, with the handles usually bearing the trademark of the business and perhaps some advertising message as well. Some manufacturers of these openers have configured the handles to identify with specific businesses. For example, the handle may resemble a house, and openers having such handles would appeal to real estate companies for use as promotional items. The handles of others may resemble trucks and, of course, would appeal to trucking companies. U.S. Pat. Des. 329,184, Des. 341,307, Des. 342,008, Des. 354,214, Des. 355,108, Des. 368,010, and Des. 364,547 illustrate openers having handles configured for specific types of businesses. Most manufacturers of letter openers offer only a few shapes and, rarely, customize an opener for any customer. After all, the mold for producing any opener is costly, with much of the cost being attributable to the region of the mold in which the slitting portion of the opener is formed, that is the portion with the spear and the embedded blade, for that portion has shapes more complex than the handle portion. The present invention resides in an opener having a handle portion and a slitting or cutting portion which are formed separately and thereafter joined together. The invention also resides in the method of making the opener, that is to say, forming its handle and cutting portions separately and then joining them together.
{ "pile_set_name": "USPTO Backgrounds" }
A semiconductor integrated circuit typically includes metallized interconnection electrodes in a pattern which overlies regions of both relatively thick and thin insulating layers. Typically a thick insulating layer region comprises a relatively thick layer of silicon dioxide, called "field oxide," located on a major planar surface of the silicon body in which the circuit is integrated; whereas a thin insulating layer region typically comprises a relatively thin silicon dioxide layer, called "gate oxide," located on the same major planar surface of the silicon body in regions complementary to the thick insulating layer regions (field oxide regions). Insulated gate field effect transistors are located in these regions of gate oxide. At the boundary between the thick and thin insulating regions a step is formed having an oxide sidewall surface (stemming from the field oxide) over which the metallized interconnection electrode must go. A typical such situation, somewhat simplified for purposes of clarity only, is illustrated in FIG. 1. (The broken lines 4--4 and 5--5 are useful in defining typical cross sections.) Here a top view in perspective shows a monocrystalline silicon body 10; a thin central gate oxide layer 11; a thick oxide layer having a left-hand field oxide portion 12 and a right-hand field oxide portion 13; a first step 14 located at the boundary of the left-hand field oxide portion 13 with the gate oxide layer 11; a second step 15 located at the boundary of the gate oxide 11 with the right-hand portion 13; a first electrode 17 running along a path extending from the left-hand field oxide portion 12 to a termination at a point overlying the central gate oxide layer 11; and a second electrode 16 running along a path extending from the left-hand field oxide 12 to the right-hand field oxide 15. Thus the first electrode 16 crosses over both steps 14 and 15, whereas the second electrode 17 crosses over only the first step 14, for illustrative purposes only. For illustrative purposes, the first electrode 16 can be a gate electrode metallization which interconnects in two directions, and the second electrode 17 can be a gate electrode which interconnects in but one direction. When these electrodes 16 and 17 are fabricated by conventional deposition of metal silicide on polycrystalline silicon followed by selective masking and etching, there is an undesirable tendency for the electrodes 16 and 17 to be joined (shorted) together by a residual metal strip running along the bottom of the sidewall of the step, because of the tendency of the etching not te remove completely some portions of the electrodes at the step. It would therefore be desirable to have a method for fabricating interconnection which avoids this problem.
{ "pile_set_name": "USPTO Backgrounds" }
1. Field of the Invention The present invention relates to a leak detecting method for vessels such as pressure vessels and hermetically sealed vessels and more particularly to a leak detecting method for vessels which employs sulfur hexafluoride as a tracer gas to utilize the absorption of a carbon dioxide gas P(16)-line laser light by the tracer gas. 2. Description of the Prior Art Included among conventional leak detecting methods for vessels represented for example by pressure vessels and hermetically sealed vessels is a method so designed that as for example, a Freon gas is sealed into a vessel and a halogen detector is used to detect any leakage from the vessel. There is known another method so designed that the whole vessel is placed within a high vacuum vessel and a helium gas is sealed into the vessel, thereby detecting any leakage from the vessel by use of a helium detector. With the conventional leak detecting methods for vessels, the former prior art or the method of using a Freon gas has the disadvantage of destroying the terrestrial environment, and thus its use will be limited or inhibited in the future. Also, the method of using a helium gas is disadvantageous in that since the vessel is placed within the high vacuum vessel, a test equipment is considerably increased in size and moreover it is difficult to reduce the locations of any leaks.
{ "pile_set_name": "USPTO Backgrounds" }
The present invention relates generally to the art of forming thin films, and more particularly to improvements in the magnetron sputtering process and apparatus for forming such films which may be highly permeable materials. One technique for depositing thin films of a desired material on a substrate is diode sputtering. A target comprising the material to be deposited, is bombarded by gas ions which have been accelerated by an intense electric field. The bombardment ejects atomic sized particles of the target which settle upon the substrate surface as a thin film. This sputtering process is slow compared to other techniques and the electric voltage required to produce a diode sputtered film is relatively high. The current saturates at a low value. Disadvantages associated with the diode sputtering process have been alleviated to a large degree by the use of magnetron sputtering. As can be seen in FIG. 1A, an array of magnets 10 and 12 is positioned behind a low permeability target material 14 where the magnetron may produce a discharge of "racetrack" shape and where the magnets may be of the type disclosed in U.S. Pats. 4,162,954, 4,180,450 and 4,265,729, issued to Charles F. Morrison, Jr., which patents are incorporated herein by reference. Coupling plate 16 serves to short the magnetic fields between the two magnets at the lower portion thereof. Because of the low permeability of the target material, the magnetic lines of force 18 extend from the magnets and pass through the target material 14 and travel substantially parallel to the plane of the target surface for a certain distance. An electric field is established perpendicular to at least a portion of the magnetic field. Gas ions are accelerated by the electric field and strike target 14 causing it to eject atomic sized particles as in diode sputtering. However, the magnetic field above the target surface confines secondary electrons ejected from the target to the vicinity of the target surface and thus accelerates the rate of collisions between the secondary electrons and gas molecules of the gas plasma (generally argon). These additional collisions serve to generate additional gas ions and, hence, more gas plasma which is confined to the vicinity of the target surface. Thus, the deposition rate of magnetron sputtering over that of diode sputtering is increased by an order of magnitude. It can be seen that the looping magnetic field as indicated by lines of force 18 is necessary to trap the plasma near the surface of target 14. However, if it is desirable to sputter a high permeability material with magnetron sputtering, the looping magnetic field will be short circuited as shown in FIG. 1B. Effectively the high permeability target 24 couples all of the magnetic lines of force from one magnet to the other just as does the coupling plate 16. The lack of the looping magnetic field 18 to trap the plasma in the vicinity of the high permeability target material would reduce the magnetron sputtering to that of ordinary diode sputtering with its attendant relatively slow sputter rate due to current saturation. A number of solutions have been attempted to obtain magnetron sputtering of highly permeable materials with only limited success. In one embodiment, a very thin high permeability target is utilized so as to become saturated by the magnets and thus incapable of shunting all of the magnetic field. Unfortunately, if the targets are made thin enough such that the magnets do not shunt virtually all of the field, the targets are rapidly depleted before a film is accumulated on substantial quantities of receiving substrate. Other approaches are to utilize relatively normal target thicknesses but in conjunction with high strength magnets again serving to saturate the target material and maintain a weak magnetic field looping thereover. This generally requires at least a second set of magnets or an extremely powerful electromagnet. This works reasonably well with moderate sized targets of iron and nickel but is generally inadequate for Permalloy, Samarium cobalt, and other very high permeability materials. This approach is generally described in my co-pending patent application Ser. No. 28,434, filed Apr. 9, 1979. One further method of permitting magnetron sputtering is to reduce the strength of field required to saturate the target material. This can be accomplished by heating the target material to above its Curie point and this is discussed in U.S. Pat. No. 4,299,678 issued to Meckel on Nov. 10, 1981. However, none of the above methods lend themselves to serious industrial coating and thus most sputtering of highly permeable materials is still done by diode sputtering with its very slow rates of accumulation.
{ "pile_set_name": "USPTO Backgrounds" }
Field of the Invention The present invention relates to video coding. In particular, the present invention relates to video coding techniques associated with loop filtering and processing across slice or tile boundaries. Description of the Related Art Motion estimation is an effective inter-frame coding technique to exploit temporal redundancy in video sequences. Motion-compensated inter-frame coding has been widely used in various international video coding standards. The motion estimation adopted in various coding standards is often a block-based technique, where motion information such as coding mode and motion vector is determined for each macroblock or similar block configuration. In addition, intra-coding is also adaptively applied, where the picture is processed without reference to any other picture. The inter-predicted or intra-predicted residues are usually further processed by transformation, quantization, and entropy coding to generate a compressed video bitstream. During the encoding process, coding artifacts are introduced, particularly in the quantization process. In order to alleviate the coding artifacts, additional processing can be applied to reconstructed video to enhance picture quality in newer coding systems. The additional processing is often configured in an in-loop operation so that the encoder and the decoder may derive the same reference pictures. FIG. 1 illustrates an exemplary adaptive inter/intra video coding system incorporating in-loop filtering process. For inter-prediction, Motion Estimation (ME)/Motion Compensation (MC) 112 is used to provide prediction data based on video data from other picture or pictures. Switch 114 selects Intra Prediction 110 or inter-prediction data from ME/MC 112 and the selected prediction data is supplied to Adder 116 to form prediction errors, also called prediction residues or residues. The prediction error is then processed by Transformation (T) 118 followed by Quantization (Q) 120. The transformed and quantized residues are then coded by Entropy Encoder 122 to form a video bitstream corresponding to the compressed video data. The bitstream associated with the transform coefficients is then packed with side information such as motion, mode, and other information associated with the image unit. The side information may also be processed by entropy coding to reduce required bandwidth. Accordingly, the side information data is also provided to Entropy Encoder 122 as shown in FIG. 1 (the motion/mode paths to Entropy Encoder 122 are not shown). When the inter-prediction mode is used, a previously reconstructed reference picture or pictures have to be used to form prediction residues. Therefore, a reconstruction loop is used to generate reconstructed pictures at the encoder end. Consequently, the transformed and quantized residues are processed by Inverse Quantization (IQ) 124 and Inverse Transformation (IT) 126 to recover the processed residues. The processed residues are then added back to prediction data 136 by Reconstruction (REC) 128 to reconstruct the video data. The reconstructed video data may be stored in Reference Picture Buffer 134 and be used for prediction of other frames. As shown in FIG. 1, incoming video data undergoes a series of processing in the encoding system. The reconstructed video data from REC 128 may be subject to various impairments due to the series of processing. Accordingly, various loop processing is applied to the reconstructed video data before the reconstructed video data is used as prediction data in order to improve video quality. In the High Efficiency Video Coding (HEVC) standard being developed, Deblocking Filter (DF) 130, Sample Adaptive Offset (SAO) 131 and Adaptive Loop Filter (ALF) 132 have been developed to enhance picture quality. The Deblocking Filter (DF) 130 is applied to boundary pixels and the DF processing is dependent on the underlying pixel data and coding information associated with the corresponding blocks. There is no DF-specific side information needs to be incorporated in the video bitstream. On the other hand, the SAO and ALF processing are adaptive, where filter information such as filter parameters and filter type may be dynamically changed according to the underlying video data. Therefore, filter information associated with SAO and ALF is incorporated in the video bitstream so that a decoder can properly recover the required information. Furthermore, filter information from SAO and ALF is provided to Entropy Encoder 122 for incorporation into the bitstream. In FIG. 1, DF 130 is applied to the reconstructed video first; SAO 131 is then applied to DF-processed video; and ALF 132 is applied to SAO-processed video. However, the processing order among DF, SAO and ALF may be re-arranged. In the High Efficiency Video Coding (HEVC) video standard being developed, the loop filtering process includes DF and SAO. The coding process in HEVC is applied to each Largest Coding Unit (LCU). The LCU is adaptively partitioned into coding units using quadtree. Therefore, the LCU is also called coding tree block (CTB). In each leaf CU, DF is performed for each 8×8 block and in HEVC Test Model Version 5.0 (HM-5.0), the DF is applied to the 8×8 block boundaries. For each 8×8 block, horizontal filtering across vertical block boundaries is first applied, and then vertical filtering across horizontal block boundaries is applied. Sample Adaptive Offset (SAO) 131 is also adopted in HM-5.0, as shown in FIG. 1. SAO is regarded as a special case of filtering where the processing only applies to one pixel. To apply SAO, a picture may be divided into multiple LCU-aligned regions. Each region can select one SAO type among two Band Offset (BO) types, four Edge Offset (EO) types, and no processing (OFF). For each to-be-processed (also called to-be-filtered) pixel, BO uses the pixel intensity to classify the pixel into a band. The pixel intensity range is equally divided into 32 bands, as shown in FIG. 2. Four consecutive bands are grouped together, where the starting band is indicated by sao_band_position. An exemplary 4-band group 200 is illustrated in FIG. 2. The first band position of this 4-band group is indicated by arrow 210. In EO, pixel classification is first done to classify pixels into different groups (also called categories or classes). The pixel classification for each pixel is based on a 3×3 window, as shown in FIG. 3 where four configurations corresponding to 0°, 90°, 135°, and 45° are used for classification. Upon classification of all pixels in a picture or a region, one offset is derived and transmitted for each group of pixels. In HM-5.0, SAO is applied to luma and chroma components, and each of the luma components is independently processed. Similar to BO, one offset is derived for all pixels of each category except for category 4 of EO, where Category 4 is forced to use zero offset. Table 1 below lists the EO pixel classification, where “C” denotes the pixel to be classified. TABLE 1CategoryCondition0C < two neighbors1C < one neighbor && C == one neighbor2C > one neighbor && C == one neighbor3C > two neighbors4None of the above Adaptive Loop Filtering (ALF) 132 is another in-loop filtering in HM-5.0 to enhance picture quality, as shown in FIG. 1. Multiple types of luma filter footprints and chroma filter footprints are used. The ALF operation is applied in the horizontal direction first. After horizontal ALF is performed, ALF is applied in the vertical direction. In HM-5.0, up to sixteen luma ALF filters and at most one chroma ALF filter can be used for each picture. In order to allow localization of ALF, there are two modes for luma pixels to select filters. One is a Region-based Adaptation (RA) mode, and the other is a Block-based Adaptation (BA) mode. In addition to the RA and BA for adaptation mode selection at picture level, Coding Units (CUs) larger than a threshold can be further controlled by filter usage flags to enable or disable ALF operations locally. As for the chroma components, since they are relatively flat, no local adaptation is used in HM-5.0, and the two chroma components of a picture share the same filter. In MH-5.0, an ALF filter for a region may be selected from multiple ALF filters. In addition, multiple filter footprints are used in HM-5.0. For each ALF filter, there is a set of filter coefficients associated with the filter. Therefore, the ALF information comprises identification for the selected ALF filter, the filter footprint and filter coefficients. As shown in FIG. 1, DF 130 is applied to reconstructed pixels from REC 128. SAO 131 is then applied to DF-processed pixels and ALF 132 is applied to SAO-processed pixels. While the processing sequence illustrated in FIG. 1 is DF, SAO and ALF, other processing sequence may also be used. For example, SAO may be applied to reconstructed pixels from REC 128, DF-processed reconstructed pixels (i.e., DF applied to the reconstructed pixels), ALF-processed reconstructed pixels (i.e., ALF applied to reconstructed pixels), both DF-processed and ALF-processed pixels (i.e., DF applied to the reconstructed pixels and ALF applied to the DF-processed reconstructed pixels) or both ALF-processed and DF-processed pixels (i.e., ALF applied to the reconstructed pixels and DF applied to the ALF-processed reconstructed pixels). For convenience, the “processed-reconstructed pixels” may refer to any type of the processed pixels mentioned above during SAO processing. The “processed-reconstructed pixels” also includes the reconstructed pixels from REC 128. In this case, it can be considered that a null processing is applied to the reconstructed pixels from REC 128. Similarly, the “processed-reconstructed pixels” may also refer to various types of the processed pixels by DF, SAO, both DF and SAO or both SAO and DF during ALF processing. Again, for ALF processing, the “processed-reconstructed pixels” also includes the reconstructed pixels from REC 128. To reduce side-information associated with SAO processing, SAO information of a current LCU can reuse the SAO information of a neighboring LCU above or to the left of the current LCU. The SAO information sharing is indicated by merge syntax. In HM-8.0, SAO syntax consists of sao_merge_left_flag, sao_merge_up_flag, sao_type_idx_luma, sao_type_index_chroma, sao_eo_class_luma, sao_eo_class_chroma, sao_band_position, sao_offset_abs, and sao_offset_sign, as shown in Table 2. Syntax sao_merge_left_flag indicates whether the current LCU reuses the SAO parameters of the left LCU. Syntax sao_merge_up_flag indicates whether the current LCU reuses the SAO parameters of the upper LCU. Syntax sao_type_idx represents the selected SAO type (sao_type_idx_luma and sao_type_idx_chroma for luma component and chroma component respectively). Syntax sao_offset_abs represents the offset magnitude and syntax sao_offset_sign represents the offset sign. Syntax cIdx indicates one of three color components. Similar mechanism can also be used to allow neighboring blocks to share the same ALF information. TABLE 2Descriptorsao( rx, ry ){if( rx > 0 ) {leftCtbInSliceSeg = CtbAddrInSliceSeg > 0leftCtbInTile = ( TileId[ CtbAddrInTS ]= = TileId[ CtbAddrRStoTS[ CtbAddrInRS − 1 ] ] )if( leftCtbInSliceSeg && leftCtbInTile )sao_merge_left_flagae(v)}if( ry > 0 && !sao_merge_left_flag ) {upCtbInSliceSeg = ( CtbAddrInRS − PicWidthInCtbsY ) >= slice_segment_addressupCtbInTile = ( TileId[ CtbAddrInTS ] = =TileId[ CtbAddrRStoTS[ CtbAddrInRS − PicWidthInCtbsY ] ] )if( upCtbInSliceSeg && upCtbInTile )sao_merge_up_flagae(v)}if( !sao_merge_up_flag && !sao_merge_left_flag ) {for( cIdx = 0; cIdx < 3; cIdx++ ) {if( ( slice_sao_luma_flag && cIdx = = 0 ) | |( slice_sao_chroma_flag && cIdx > 0 ) ) {if( cIdx = = 0 )sao_type_idx_lumaae(v)else if( cIdx = = 1 )sao_type_idx_chromaae(v)if( SaoTypeIdx[ cIdx ][ rx ][ ry ] != 0 ) {for( i = 0; i < 4; i++ )sao_offset_abs[ cIdx ][ rx][ ry ][ i ]ae(v)if( SaoTypeIdx[ cIdx ][ rx ][ ry ] = = 1 ) {for( i = 0; i < 4; i++ )if( sao_offset_abs[ cIdx ][ rx ][ ry ][ i ] != 0 )sao_offset_sign[ cIdx ][ rx ][ ry ][ i ]ae(v)sao_band_position[ cIdx ][ rx ][ry ]ae(v)} else {if( cIdx = = 0 )sao_eo_class_lumaae(v)if( cIdx = = 1 )sao_eo_class_chromaae(v)}}}}}} The LCUs in a picture can be partitioned into slices, where each slice consists of multiple horizontally consecutive LCUs. In HM-5.0, another image unit structure, named tile, is introduced, where a picture is partitioned into multiple tiles. For example, a picture may be divided into M tiles horizontally and N tiles vertically, where M and N are integers greater than 0. Each tile consists of multiple LCUs. Within each tile, the processing sequence of the LCUs is according to the raster scan order. Within each picture, the processing sequence of the tiles is also according to the raster scan order. Tile boundaries are often aligned with LCU boundaries. In some systems, it is desirable to process the slices or tiles independently. Independent slice/tile processing will allow parallel processing of multiple slices or tiles. For CTBs or LCUs located at a left boundary or a top boundary of the slice or tile, SAO or ALF parameter sharing with a neighboring LCU above or to the left of the current LCU implies data dependency on an LCU from another slice or tile. Therefore, it is desirable to develop SAO or ALF processing that enables slice/tile independent processing.
{ "pile_set_name": "USPTO Backgrounds" }
The present invention generally relates to semiconductor devices (here, the invention chiefly relates to semiconductor memory devices, and a semi-conductor memory device will hereinafter be described), and particularly to the redundancy technique for relieving defects by replacing defective memory cells by spare memory cells. High-density integration of semiconductor memory has been highly advanced up to the mass production of dynamic random access memories (DRAM) of 256 mega bits. This high-density integration advancement involves extreme reduction of element size and great increase of the number of elements, thus causing a problem of low yield due to defects. As a counter-measure against this problem, the so-called redundancy technique is known, in which the defective memory cells are replaced, or repaired by redundant memory cells as the spare, or back-up memory cells that are previously provided on a memory chip. Efforts to enhance the efficiency of the defect repair have been made in this technical field. An example of the defect repair technology for DRAM is disclosed in JP-A-2-192100 (laid open Jul. 27, 1990) in which the decision for column-side repair performed according to a row address, and column selection lines are replaced by redundant column selection lines so that block repair can be made. This method is powerful because a large number of defective memory cells can be replaced by a small number of redundant column selection lines. FIG. 2 is a schematic diagram of a conventional logic construction for block redundancy. Here, a memory cell group of two regions including defects is replaced by a redundant memory cell group. A memory cell array NMCA and a redundant cell array RMCA are provided and controlled by a repair decision circuit YRC. The memory cell array NMCA has memory cells provided at the intersections of N word lines WLs and M data lines DLs, and the memory cells are selected by a row decoder XDEC and a column decoder YDEC. The redundant cell array RMCA has redundant memory cells provided at the intersections of N word lines WLs and P data lines RDLs, and the redundant memory cells are selected by the row decoder XDEC and a redundant column decoder RYD. The row decoder XDEC decodes a row address AX of n bits and selectively drives one of the 2n, or N word lines. The column decoder YDEC decodes a column address AY of m bits and selects one of the 2m, or M data lines DLs. The redundant column decoder RYD decodes p bits of the column address AY, and selects one of the 2P, or P redundant data lines RDLs. A repair decision result RYH from the repair decision circuit YRC controls the column decoder YDEC and the redundant column decoder RYD. If the repair decision result RYH is xe2x80x980xe2x80x99,the column decoder YDEC is activated to select memory cells within the memory cell array NMCA. If the repair decision result RYH is xe2x80x981xe2x80x99, the redundant column decoder RYD is activated to select redundant memory cells within the redundant cell array RMCA. Thus, the memory cell group of defects DF1, DF2 can be replaced by a redundant memory cell group. A unit of memory cells to be replaced is the area selected by Q word lines and P data lines. The repair decision circuit YRC is formed by two row address comparators AXC, two column address comparators AYC, two dual-input AND gates AND2, and an dual-input OR gate OR2. A pair of one row address comparator AXC and one column address comparator AYC stores one region to be repaired, or replaced. Each row address comparator AXC includes address storage means for storing a repair address of (nxe2x88x92q) bits, and compares it with the (nxe2x88x92q) bits of the row address AX. Each column address comparator AYC includes address storage means for storing a repair address of (mxe2x88x92p) bits, and compares it with the (mxe2x88x92p) bits of the column address AY. The dual-input AND gates AND2 take logic products of coincidence decision results XHC1, XHC2 from the row address comparators AXC and coincidence decision results YH1, YH2 from the column address comparators AYC to produce decision results HC1, HC2 for the two, first and second replacements. The dual-input OR gate OR2 takes a logical sum of these decision results to produce the repair decision result RYH. Since the repair decision circuit is constructed as above, defects at separate column addresses can be repaired according to the row addresses, or replaced by redundant memory cells on the same redundant data line. In the column block redundancy shown in FIG. 2, the row addresses in the first replacement must be different from those in the second replacement. In other words, the repair row address stored in one of the two row address comparators AXC must be different from that of the other. If the same row address were stored in the two comparators, the replacing regions RPD would be one region, or the replaced regions would compete with each other for acquiring the one region irrespective of whether the column addresses of the replaced regions RPO are different or not. Therefore, even though two repair addresses can be stored, it is impossible to repair the case in which two defects occur in different-column-address regions but in the same-row-address regions each of which is selected by Q word lines and P data lines. In order to reduce the probability of that case in which both defects cannot be relieved because the replaced regions RPO compete with each other about taking one replacing region RPD, it can be considered to decrease the number, Q of word lines that are one replacement unit. However, if the number Q is decreased, it is necessary to increase the number of bits, (nxe2x88x92q) of the row address that the row address comparator AXC compares with, so that the circuit scale of the row address comparator becomes large. Accordingly, it is desired to contrive a method of effectively repairing a plurality of defects at the same time. That is, it is an object of the invention to provide a semiconductor memory device having a redundancy circuit capable of effectively repairing defects by use of small-circuit-scale address comparators that compare with a smaller number of bits, and by controlling the replacement operation so that the competition between the replaced regions can be avoided. According to one aspect of the invention, there is provided a semiconductor memory device having a plurality of word lines, a plurality of bit lines arranged to intersect the plurality of word lines, a large number of memory cells arranged at necessary intersections between the plurality of word lines and the plurality of bit lines, a plurality of spare bit lines arranged to intersect the plurality of word lines, a plurality of spare memory cells arranged at necessary intersections between the plurality of word lines and the plurality of spare bit lines, and a redundancy circuit for replacing memory cell groups, including defects, of the large number of memory cells by spare memory cell groups of the spare memory cells, wherein the redundancy circuit has functions to control a first replacement to be made by a first replacing unit, and a second replacement to be made by a second replacing unit that is smaller than the first replacing unit, and to give the second replacement priority when the first and second replacements compete with each other about taking the replacing spare memory cell groups. In other words, the redundancy circuit controls the first replacement to be made by the first replacing unit and the second replacement to be made by the second replacing unit that is smaller than the first replacing unit, and it includes a first decision circuit for deciding about at least a first part of an address provided to select the large number of memory cells, a second address decision circuit for deciding about a second part of the address, and a third address decision circuit for deciding about at least a third part of the address except the second part, whereby when the second address decision circuit produces a miss, the second replacement is not performed but the first replacement is made according to the output from the first address decision circuit, and when the second address decision circuit produces a hit, the first replacement is not performed but the second replacement is made according to the output from the third address decision circuit.
{ "pile_set_name": "USPTO Backgrounds" }
1. Technical Field The present invention relates generally to high-performance, fault-tolerant HTTP, streaming media and applications delivery in a content delivery network (CDN). 2. Description of the Related Art It is well-known to deliver HTTP and streaming media using a content delivery network (CDN). A CDN is a self-organizing network of geographically distributed content delivery nodes that are arranged for efficient delivery of digital content (e.g., Web content, streaming media and applications) on behalf of third party content providers. A request from a requesting end user for given content is directed to a “best” replica, where “best” usually means that the item is served to the client quickly compared to the time it would take to fetch it from the content provider origin server. An entity that provides a CDN is sometimes referred to as a content delivery network service provider or CDNSP. Typically, a CDN is implemented as a combination of a content delivery infrastructure, a request-routing mechanism, and a distribution infrastructure. The content delivery infrastructure usually comprises a set of “surrogate” origin servers that are located at strategic locations (e.g., Internet Points of Presence, access points, and the like) for delivering copies of content to requesting end users. The request-routing mechanism allocates servers in the content delivery infrastructure to requesting clients in a way that, for web content delivery, minimizes a given client's response time and, for streaming media delivery, provides for the highest quality. The distribution infrastructure consists of on-demand or push-based mechanisms that move content from the origin server to the surrogates. An effective CDN serves frequently-accessed content from a surrogate that is optimal for a given requesting client. In a typical CDN, a single service provider operates the request-routers, the surrogates, and the content distributors. In addition, that service provider establishes business relationships with content publishers and acts on behalf of their origin server sites to provide a distributed delivery system. A well-known commercial CDN service that provides web content and media streaming is provided by Akamai Technologies, Inc. of Cambridge, Mass. CDNSPs may use content modification to tag content provider content for delivery. Content modification enables a content provider to take direct control over request-routing without the need for specific switching devices or directory services between the requesting clients and the origin server. Typically, content objects are made up of a basic structure that includes references to additional, embedded content objects. Most web pages, for example, consist of an HTML document that contains plain text together with some embedded objects, such as .gif or jpg images. The embedded objects are referenced using embedded HTML directives, e.g., Uniform Resource Identifiers (URIs). A similar scheme is used for some types of streaming content which, for example, may be embedded within an SMIL document. Embedded HTML or SMIL directives tell the client to fetch embedded objects from the origin server. Using a CDN content modification scheme, a content provider can modify references to embedded objects so that the client is told to fetch an embedded object from the best surrogate (instead of from the origin server). In operation, when a client makes a request for an object that is being served from the CDN, an optimal or “best” edge-based content server is identified. The client browser then makes a request for the content from that server. When the requested object is not available from the identified server, the object may be retrieved from another CDN content server or, failing that, from the origin server. In some CDNs, such as Akamai FreeFlow® content delivery service, data about the content provider's (CP's) objects, or so-called “metadata,” is often directly encoded “in-URL,” namely in the HTML or SMIL directives that are modified during the content modification process. More specifically, metadata is the set of all control options and parameters that determine how a CDN content server will handle a request for an object. Such metadata may include, for example, a CP code or other internal tracking number used, for example, to facilitate billing, coherence information (e.g., TTL or fingerprint) about how CDN servers should cache the object and maintain its freshness, a unique serial number value that may be used for load balancing, access control data, a hostname identifying the origin server where a copy of the object may be located, and other feature-specific metadata. By including object metadata directly in the HTML or SMIL directives, content providers may set up their metadata as part of the publication process, i.e., without requiring changes in their web server or involving network operations personnel. The “in-URL” embedding technique ensures that any modified URL pointing to the CDN has, in a self-contained way, the information needed to serve the object. On the other hand, the modified URL generated by this process is often long and complex. In addition, some content providers may only have a need to specify site-wide or global metadata specifications. Thus, there remains a need to provide a framework that allows for both a simple method of creating modified URLs for sites with simple global metadata specifications, while allowing arbitrary complexity for sites with arbitrarily complex metadata needs.
{ "pile_set_name": "USPTO Backgrounds" }
1. Field of the Invention The present invention relates generally to a method and an apparatus for soldering surface mounted electronic components to and desoldering them from metallic contacts on substrates such as printed circuit boards. 2. Summary of the Invention and Description of the Prior Art A convection heat system for metallic solder attachment and detachment of components surface-mounted on printed circuit boards is disclosed. For soldering, the system conveniently, rapidly and precisely positions the workpiece component, the printed circuit board and a source of convection heat relatively to each other through either moving the printed circuit board using X-Y table or through moving the vacuum-held workpiece with the source of convection heat to a desired relation with the printed circuit board using an arm system that moves in X, Y and Z directions independently of the X-Y table. The source of convection heat used for melting the solder is hot gas which is forced to flow down the terminus of the arm inside a nozzle conduit that encompasses the component work area where the hot gas deflecting inwardly from the nozzle conduit walls moves rapidly past the area to be heated and then outward from the printed circuit board working area. At the same time that the hot gas flow is melting solder, cooling air is forced down a tube inside the nozzle conduit onto the top central portion of the workpiece component surface. The tube used for passage of the cooling air may advantageously be used also for the vacuum holding of the workpiece since the two operations take place at different times. A thermal sensor in the immediate vicinity of the solder joints effectively measures the temperature of the solder so that the heating process can be controlled or terminated. Hot gas has been used in a number of existing systems to melt solder contacts between surface mounted components and printed circuit boards so as to release or secure these contacts. One such system was disclosed in the U.S. Pat. No. 4,426,571, where the main problem is that the hot air flow is not restricted to the immediate vicinity of the workpiece, and overheating of components adjacent to the workpiece occurs. Large amount of hot air required in this case often causes damage to the printed circuit board. In the present invention hot gas is restricted to the immediate vicinity of the workpiece by means of a nozzle conduit. Another problem encountered in practical applications of the method disclosed in the U.S. Pat. No. 4,426,571 is damage of circuitry printed on the printed circuit board. This damage may easily occur when the workpiece is picked up by hand or a hand held tool while the solder contacts connecting the workpiece to the printed circuitry are not completely melted. The present invention provides vacuum pick-up so that no excessive force may be applied to lift the workpiece off the printed circuit board. The idea of restricting hot gas to the immediate vicinity of the workpiece has been previously disclosed in the U.S. Pat. No. 4,552,300, where hot gas fills the cavity formed between the nozzle walls and the printed circuit board, and the workpiece is located in the cavity. This method proved to be inefficient in transferring heat from hot gas to solder contacts. In the present invention the heat exchnge occurs between hot gas rapidly moving directly past the solder contacts into the area outside the nozzle conduit. Another problem that often occurs when soldering or desoldering contacts between the workpiece and the printed circuit board is overheating and/or thermal shock damage of the workpiece itself. The present invention provides means for cooling the workpiece and shielding it from direct incidence of hot gas to avoid overheating and thermal shock damage of the workpiece. In systems such as CRAFT-100 by PACE Inc. and systems manufactured by AIRVAC little or no visual and physical access is provided to the workpiece during positioning of the workpiece on the printed circuit board. In the present invention convenient visual and physical access is provided to the workpiece at any time during positioning by allowing the nozzle conduit to move away from the workpiece during positioning and completely enclose the workpiece during soldering. In existing systems, relative positioning of the workpiece and the printed circuit board is accomplished by moving the printed circuit board so that the area of interest is precisely underneath the component and the source of heat. Moving the printed circuit board, which is often many times larger than the component, in precise manner is often inconvenient. The present invention provides means for moving the workpiece held by vacuum at the end of an arm mechanism together with the source of heat as well as means for moving the printed circuit board to position the workpiece, the printed circuit board and the source of heat relatively to each other. It is, thus, a principle object of the present invention to provide an apparatus for soldering and desoldering of contacts between surface mounted components and printed circuit boards that is more convenient and efficient than existing apparatuses. It is an important object of this invention to provide an apparatus as described in which an efficient transfer of heat from hot gas to solder contacts is achieved by creating a turbulent rapid flow of hot gas directly past the solder contacts. It is a further object of this invention to provide an apparatus as described where the printed circuit board and the components adjacent to the workpiece are protected from overheating by restricting the application of hot gas to the immediate vicinity of the workpiece. Another important object is to provide an apparatus as described in which the workpiece itself is protected from overheating and thermal shock by providing means of cooling the workpiece and shielding it from direct incidence of hot gas during the process of soldering and desoldering. It is a further object of this invention to provide an apparatus as described in which the hot gas is delivered to the solder contacts at high flow velocity without significant impact, eliminating splattering of the solder, by distributing uniformly the hot gas flow along and deflecting the flow from the walls of the nozzle conduit used to restrict hot gas application. It is an important object of this invention to provide a method and an apparatus to position precisely and conveniently the surface mounted components, the printed circuit board and the hot gas flow relatively to each other. It is yet another object to provide an apparatus as described above in which there is an easy and convenient visual and physical access to the workpiece and its location on the printed circuit board by allowing the nozzle conduit to move relatively to the vacuum-held workpiece.
{ "pile_set_name": "USPTO Backgrounds" }
In a conventional wireless communication system including a pair of transceivers communicating with one another over a wireless communication channel, there are typically a number of different data transmission rates available at which to transmit data. Generally, the higher the data rate, the more susceptible the system is to errors. Under certain circumstances, it is necessary to adapt the system to higher or lower data transmission rates based, at least in part, on environmental conditions. For example, noise on the communication channel, transceiver impairments, etc., may necessitate operation of the system at a lower data transmission rate. The Institute of Electrical and Electronics Engineers (IEEE) 802.11 standard addresses medium access control over a wireless local area network (WLAN). The IEEE 802.11 standard is set forth in the document IEEE Std. 802.11, entitled Supplement to IEEE Standard for Information Technology—Telecommunications and Information Exchange Between Systems—Local Metropolitan Area Networks—Specific Requirements—Part 11: Wireless LAN Medium Access Control (MAC) and Physical Layer (PHY) Specifications, 1999 Edition, which is incorporated herein by reference. Additional extensions relating to the 802.11 standard, including IEEE Std. 802.11a, entitled High Speed Physical Layer in the 5 GHz Band, February 2000, and IEEE Std. 802.11g, entitled Further Higher Data Rate Extension in the 2.4 GHz Band, September 2000, are also incorporated herein by reference. Rate adaptation in a wireless communication system operating in accordance with the 802.11 standard generally takes place in the transmitter at the MAC level. Known rate adaptation schemes typically rely on information acquired through acknowledgment messages received after each correctly transmitted data packet. An acknowledgment message indicates a correctly received packet, while an absence of an acknowledgment message is generally interpreted as an error. A determination as to whether to change the data rate in the transmitter can be made in response to the number of consecutive acknowledgments that are received. After a certain number of correctly received data packets, the transmitter typically attempts to switch to a higher data transmission rate. Similarly, after a certain number of consecutive errors, the transmitter attempts to switch to a lower data transmission rate. This conventional rate-switching methodology, which is based on received acknowledgments, has the advantage of simplicity. However, it often changes the data transmission rate of the transmitter to a value that is either too high or too low, thus undesirably impacting the throughput of the system. For example, switching to a lower data rate when, in fact, a higher rate can be supported by the system results in a significant throughput degradation. The same is true when switching to a higher data rate than the system can support, thus resulting in a high packet error rate (PER), bit error rate (BER), or frame error rate (FER). It would be desirable, therefore, to be able to obtain an accurate estimate of the signal quality of a received signal for, among other applications, controlling the data transmission rate in a wireless communication system, which addresses the above-mentioned problems exhibited in conventional wireless communication systems.
{ "pile_set_name": "USPTO Backgrounds" }
1. Technical Field The present invention relates to a method of estimating a channel, and more particularly, to an apparatus and a method of estimating a channel in consideration of a residual synchronization offset in which a phase generated by the residual synchronization offset in proportion to a difference between a channel estimation time and a channel compensation time is reflected in a channel estimated value in the case in which an orthogonal frequency division multiplexing (OFDM) symbol estimating the channel and an OFDM symbol compensating for the channel are different from each other. 2. Description of the Related Art In an orthogonal frequency division multiplexing (OFDM) transmission scheme, which is one of schemes of transmitting a plurality of carriers in which several carriers are used, input data are carried on a plurality of sub-carriers having orthogonality and are transmitted in parallel with each other. In the OFDM transmission scheme, a transmission period in each channel is increased by the number of carriers. In this case, frequency selective channel characteristics appearing by using a wide band at the time of transmitting high speed data are approximated to frequency non-selective channel characteristics by a narrowed channel. Therefore, in the OFDM transmission scheme, distortion by a channel may be compensated for by only an equalizer of a single sample simpler than a single carrier system, such that the OFDM transmission scheme has been widely used in a high speed data transmission system of several fields such as multimedia data transmission, and the like. In a system using the OFDM transmission scheme, a channel is estimated using a pilot sub-carrier. As a scheme of estimating a channel using a pilot sub-carrier, there are several schemes such as a linear minimum mean square error (LMMSE) scheme, which is a scheme of estimating a channel using statistical characteristics of the channel, a maximum likelihood (ML) scheme, which is a scheme of calculating an estimated value having a maximum likelihood value, a general linear and polynomial interpolation scheme, and the like. The channel is estimated by one of these several schemes of estimating a channel and is then compensated for based on the estimated channel value. However, in a scheme of estimating a channel according to the related art, an influence of a difference between an OFDM symbol estimating the channel and an OFDM symbol compensating for the channel and a residual synchronization offset generated by the difference is not considered in estimating the channel and compensating for the channel. That is, in the case in which the OFDM symbol estimating the channel and an OFDM symbol compensating for the channel are different from each other, channel estimating performance is deteriorated due to the residual synchronization offset.
{ "pile_set_name": "USPTO Backgrounds" }
1. Field of the Invention The present invention generally relates to a keypad, and in particular, to a keypad and a keypad assembly for a portable communication terminal. 2. Description of the Related Art Portable communication terminals are evolving into complex multi-functional devices capable of providing various functions. Thus, a keypad of the portable communication terminal has to support characters and special characters in addition to numbers, and their size has to be fit accordingly due to the characteristic of a portable device. A general qwerty-type keypad can be mounted in a portable communication terminal to input characters, but its large volume may become an issue due. These days, the keypad mounted in the portable communication terminal needs to provide a plurality of characters, numbers and special characters in a single key button. A general keypad mounted in the portable terminal usually serves as a number input means for selecting a phone number. A character coining method in which characters, numbers, and special characters on a display device are selected with a single key may be used. In addition, the portable communication terminal may include a touch screen as an input means. However, in a touch screen type input device, a user cannot feel the activation of a key during operation. If the user cannot feel the press of a key or cannot accurately press the desired key, the user may press adjacent another key, causing inaccurate manipulation of the keypad. Moreover, in a keypad supporting a plurality of characters with a single key, the user has to select the desired character from small-size characters printed in each key. As a result, a user who is not accustomed to use this type of keypad has to spend much time and effort in learning to manipulate the keypad accurately.
{ "pile_set_name": "USPTO Backgrounds" }
1. Field of the Invention The invention relates to the distribution of television signals. 2. Description of the Prior Art Television signals are conventionally conveyed by coaxial cables, and thus in unbalanced form, but in certain circumstances it is beneficial to be able to distribute television signals over a twisted pair of conductors, i.e. in balanced form, especially if a twisted pair of conductors is available in a cable otherwise used to distribute telephone signals or local area network signals. The invention aims to enable distribution using a twisted pair of conductors whose length may vary within a wide range.
{ "pile_set_name": "USPTO Backgrounds" }
When a procedure involving the percutaneous insertion of an instrument such as a catheter into a blood vessel is carried out for medical treatment, examination or diagnosis, bleeding at the puncture site must be stopped following subsequent withdrawal and removal of the catheter. Hemostatic devices which are attached by being wrapped around the portion of an arm or leg where the puncture site is located and compress the puncture site where bleeding is to be stopped are already known in the prior art (e.g., JP 3031486 U). In such prior-art hemostatic devices, the pressure applied to the site where bleeding is to be stopped is directed in a substantially vertically downward direction, and the hemostatic effect in this pressing direction is inadequate. Therefore, complete hemostatis sometimes does not occur or takes a long time to achieve. Moreover, prior-art hemostatic devices apply pressure not only to the puncture site where bleeding is to be stopped, but to the surrounding area as well. Hence, other tissues are also compressed, including other blood vessels and nerves, sometimes resulting in numbness and poor blood circulation. To keep this from happening, a health care practitioner such as a physician or nurse must lower the compressive force over time by carrying out manual operations to reduce the balloon pressure or loosen the band, which is inefficient and inconvenient. In addition, when using such prior-art hemostatic devices, the health care practitioner visually sights the balloon into place over the puncture site. It has been pointed out that this makes the balloon difficult to position properly. In fact, due to poor positioning of the balloon, a hematoma may form or blood leakage may occur because of the inability to stop bleeding.
{ "pile_set_name": "USPTO Backgrounds" }
Semiconductor devices are used in a variety of electronic applications, such as personal computers, cell phones, digital cameras, and other electronic equipment. Semiconductor devices are typically fabricated by sequentially depositing insulating or dielectric layers, conductive layers, and semiconductor layers of material over a semiconductor substrate, and patterning the various material layers using lithography to form circuit components and elements thereon. Many integrated circuits are typically manufactured on a single semiconductor wafer. The dies of the wafer may be processed and packaged at the wafer level, and various technologies have been developed for wafer level packaging.
{ "pile_set_name": "USPTO Backgrounds" }
The present application relates generally to registering multiple datasets with each other. In an exemplary embodiment, the present invention relates to registering differing ophthalmic datasets from different measurements (such as a wavefront measurement and a corneal topography map) of an eye. Known laser eye procedures generally employ an ultraviolet or infrared laser to remove a microscopic layer of stromal tissue from the cornea of the eye to alter the refractive characteristics of the eye. The laser removes a selected shape of the corneal tissue, often to correct refractive errors of the eye. Ultraviolet laser ablation results in photo-decomposition of the corneal tissue, but generally does not cause significant thermal damage to adjacent and underlying tissues of the eye. The irradiated molecules are broken into smaller volatile fragments photochemically, directly breaking the intermolecular bonds. Laser ablation procedures can remove the targeted stroma of the cornea to change the cornea's contour for varying purposes, such as for correcting myopia, hyperopia, astigmatism, and the like. Control over the distribution of ablation energy across the cornea may be provided by a variety of systems and methods, including the use of ablatable masks, fixed and moveable apertures, controlled scanning systems, eye movement tracking mechanisms, and the like. In known systems, the laser beam often comprises a series of discrete pulses of laser light energy, with the total shape and amount of tissue removed being determined by the shape, size, location, and/or number of a pattern of laser energy pulses impinging on the cornea. A variety of algorithms may be used to calculate the pattern of laser pulses used to reshape the cornea so as to correct a refractive error of the eye. Known systems make use of a variety of forms of lasers and/or laser energy to effect the correction, including infrared lasers, ultraviolet lasers, femtosecond lasers, wavelength multiplied solid-state lasers, and the like. Alternative vision correction techniques make use of radial incisions in the cornea, intraocular lenses, removable corneal support structures, thermal shaping, and the like. Known corneal correction treatment methods have generally been successful in correcting standard vision errors, such as myopia, hyperopia, astigmatism, and the like. However, as with all successes, still further improvements would be desirable. Toward that end, wavefront measurement instruments are now available to measure the refractive characteristics of a particular patient's eye. One promising wavefront measurement system is the iDESIGN ADVANCED WAVESCAN STUDIO System, which includes a Hartmann-Shack wavefront sensor assembly that may quantify higher-order aberrations throughout the entire optical system, including first and second-order sphero-cylindrical errors and third through sixth-order aberrations caused by coma and spherical aberrations. With advanced algorithms for measuring and applying the wavefront correction (e.g. Fourier or zonal), even higher spatial frequency structures can be corrected, providing that adequate registration can be maintained between the intended correction and its application in a practical system. The wavefront measurement of the eye creates a high order aberration map that permits assessment of aberrations throughout the optical pathway of the eye, e.g., both internal aberrations and aberrations on the corneal surface. Thereafter, the wavefront aberration information may be saved and thereafter input into a computer system to compute a custom ablation pattern to correct the aberrations in the patient's eye. A variety of alternative wavefront or other aberration measurement systems may also be available Customized refractive corrections of the eye may take a variety of different forms. For example, lenses may be implanted into the eye, with the lenses being customized to correct refractive errors of a particular patient. By customizing an ablation pattern or other refractive prescription based on wavefront measurements, it may be possible to correct minor refractive errors so as to reliably and repeatably provide visual acuities better than 20/20. Alternatively, it may be desirable to correct aberrations of the eye that reduce visual acuity, even where the corrected acuity remains less than 20/20. While wavefront measurement systems have been highly successful, improvements are still possible. For example, in some instances it may be desirable to concurrently diagnose the refractive errors of the eye using two or more different optical measurements so as to provide a better diagnosis (and treatment) of the refractive errors in the optical tissues of the eye. For example, the iDESIGN ADVANCED WAVESCAN STUDIO System includes both a wavefront aberrometer and corneal topographer. In order to fully take advantage of two different data sources for corneal treatment planning, however, it will generally be desirable to combine the data from the two optical measurement instruments. Consequently, what is needed are methods, systems and software for registering datasets from separate optical measurement devices. Multi-modal diagnostic instruments are being developed that acquire data from different measurements of the eye. For example, a multi-modal diagnostic instrument may include, for example, wavefront aberrometry and corneal topography (CT), optical coherence topography (OCT) and wavefront (WF), optical coherence topography and topography, pachymetry and wavefront, and so forth. The different measurements taken by these multi-modal diagnostic instruments may be more useful if are be registered together. Often it is difficult to acquire the images at exactly the same time, which requires synchronized cameras. Accordingly, it would be useful to provide systems and methods that allow image data to be registered that was taken by different devices, or the same device, at different times.
{ "pile_set_name": "USPTO Backgrounds" }
Increased miniaturization of the circuitry forming a wireless communication device has greatly increased the portability of such a device and has permitted carriage thereof by a user. Carrying devices that enable a user to carry a radiotelephone on a belt or other article of clothing are known in the art. These carrying devices typically include a portion for holding the radiotelephone and a clip portion attached to the carrying portion for attachment to a user's article of clothing. Once the carrying device is clipped to the article of clothing, the radiotelephone remains stationary in its angular orientation, typically vertical. This vertical orientation can cause discomfort when the user sits down. For example, the bottom of the carrying case can cause discomfort to the user's legs or the top of the carrying case or antenna of the radiotelephone can cause discomfort or pain to the user's side or stomach. Firearm holsters that can be adjusted while worn on a user's article of clothing are known. U.S. Pat. No. 3,915,361 ('361 patent) describes a holster with a case for holding a hand-gun and an adjustable mounting clip for attachment to a belt or waistband of the user's clothing. The mounting clip or paddle is rotatable about its point of attachment to the case. The orientation of the paddle relative to the case can be adjusted to allow the holster to be worn either conventionally or in a cross-draw mode. The rotatable connection includes a metal ratchet wheel that couples to a grooved recess in a bracket that is integrally formed with the paddle. The radius of the ratchet wheel is substantially the same as the radius of the recess. Tightening of a screw against the face of the bracket binds the serrated outer edges of the ratchet wheel and the recess to hold the paddle in a fixed position relative to the case. When the user desires to change the orientation of the paddle relative to the case, he or she must loosen the screw, rotate the paddle and then tighten the screw to maintain the new orientation of the paddle relative to the case. This process is timely and inconvenient because it requires the user to obtain a screwdriver or other tool before any adjustments can be made. Also, because the paddle is adjusted instead of the case, attachment to the user's belt or clothing can be compromised as the angle of orientation relative to the case increases. U.S. Pat. No. 4,504,001 describes a swivel connected belt holster wherein the rotatable connection between the belt loop assembly and the holster includes two rigid plates. One of the rigid plates includes a number of bosses in a circular array. The other plate includes the same number of mating recesses in a circular array. When the holster is in use, the bosses of one plate engage mating recesses in the other plate and the plates are held in place by a position locking screw to provide rigid angular positioning of the holster body with respect to the belt loop assembly. As with the invention of the '361 patent, adjustment of the holster with respect to the belt loop assembly requires the timely and inconvenient steps of obtaining a tool, loosening a screw, adjusting the position and tightening a screw. Accordingly, there is a need for a mountable carrying device for a portable radiotelephone whereby the position of the radiotelephone with respect to the mounting mechanism can be easily adjusted and locked into position without using a tool and while being worn on the user's belt or other article of clothing.
{ "pile_set_name": "USPTO Backgrounds" }
Image sensors, such as CCD sensors and CMOS sensors, are widely known as image sensors used for digital still cameras, digital video cameras, mobile phone cameras, endoscope cameras, etc. Currently, an image sensor having a layered structure, wherein a plurality of pixel electrodes are two-dimensional arranged on a substrate having a read-out circuit, etc., formed thereon, and an organic layer including at least a photoelectric conversion layer, and a counter electrode are formed in this order on the pixel electrodes, has been proposed (U.S. Pat. No. 7,920,189 (hereinafter, Patent Document 1)). In the image sensor having a layered structure disclosed in Patent Document 1, the photoelectric conversion layer may be a continuous layer shared by all the pixel portions, or may be divided into parts corresponding to the individual pixel portions. On the other hand, the image sensor provided with the continuous photoelectric conversion layer that is shared by the plurality of pixel electrodes has such problems that residual charges remain between the pixel electrodes and form a residual image of image information, and that the thickness of the photoelectric conversion layer is reduced at end portions of the pixel electrodes and electric field concentration occurs, resulting in leakage current, etc. In order to solve the above described problems, Japanese Unexamined Patent Publication Nos. 2008-177287 and 2009-259978 (hereinafter, Patent Documents 2 and 3, respectively), for example, propose providing the pixel electrodes having inclined end portions. Patent Document 2 teaches that it is preferable to incline the end portions of the pixel electrodes at an angle in the range from 30° to 120° relative to a planer substrate for suppressing occurrence of a residual image due to electric charges remaining between the pixel electrodes. Patent Document 3 teaches that it is preferable to incline the end portions at a predetermined inclination angle that is smaller than 90° for suppressing the decrease of the thickness of the photoelectric conversion layer at stepped portions of the end portions of the pixel electrodes and for suppressing occurrence of the electric field concentration.
{ "pile_set_name": "USPTO Backgrounds" }
1. Field of the Invention The present invention relates to an optical pickup used in an optical disk system, such as a compact disk, a video disk, etc. More particularly, the present invention relates to an optical pickup comprising a hologram element-incorporating semiconductor laser device. 2. Description of the Related Art An optical pickup comprising a semiconductor laser device is used to read out information stored in an optical disk, such as a compact disk, etc. In the optical pickup, light emitted from the semiconductor laser device is split by a diffraction grating of a hologram element into one main beam and two subbeams (tracking beams) which are brought onto an optical disk. The main beam and the subbeams are reflected on the optical disk, and each reflected beam is further split by a hologram of the hologram element into two beams, which are brought to a light receiving element or a signal processing integrated circuit with a light receiving element. Thereafter, based on an output signal from the signal processing integrated circuit, a tracking information signal, etc. used for accurately reading out signals recorded in the optical disk can be obtained. FIG. 4 is a schematic diagram showing the optical system of a conventional three-beam hologram optical pickup. This optical pickup has a semiconductor laser chip (LD) 6. Light emitted from the semiconductor laser chip 6 is split by a tracking beam generating diffraction grating 5, provided on the rear side of a hologram element (not shown), into three beams, i.e., two subbeams for tracking and one main beam for reading information signals. This light passes through a hologram 4 provided on the hologram element as zero-order light, and is then converted by a collimator lens 3 to parallel light. The parallel light is condensed by an objective lens 2 onto a disk 1. The light condensed onto the disk 1 is modulated by pits on the disk 1 and reflected from the disk 1. The reflected light from the disk 1 passes through the objective lens 2 and the collimator lens 3 in this order, and is then diffracted by the hologram 4 and introduced into a five-way split photodiode 7 as first-order diffracted light. This five-way split photodiode 7 has five optical detectors D1 to D5. The five-way split photodiode 7 has a rectangular region which is illuminated by light. The region is divided into three equal parts which are strip regions extending in a longitudinal direction. Two opposite regions are first and fifth optical detectors D1 and D5. A middle strip region is further divided into two equal parts in a transverse direction. One of the two regions is a fourth optical detector D4. The other region is further divided into two parts in a longitudinal direction, which are second and third optical detectors D2 and D3. The hologram 4 has two regions 4a and 4b which have different grating pitches. The main beam of reflected light entering the region 4a is condensed onto the splitting line between the second optical detector D2 and the third optical detector D3 of the five-way split photodiode 7. The main beam of reflected light entering the region 4b is condensed onto the fourth optical detector D4. Further, the two subbeams of reflected light entering the region 4a are condensed onto the opposite first and fifth optical detectors D1 and D5, so that two beam spots are formed on each of the optical detectors D1 and D5. As described above, the beam spots of reflected light condensed on the optical detectors D1 to D5 of the five-way split photodiode 7 vary depending on the focusing conditions of the light brought onto the disk 1 as shown in FIGS. 5A to 5C. FIG. 5A shows spots when light is focused beyond the optical disk 1. FIG. 5B shows spots when light is properly focused on the optical disk 1. FIG. 5C shows spots when light is focused before the disk 1. The outputs of the optical detectors D1 to D5 of the five-way split photodiodes 7 are represented by S1, S2, S3, S4 and S5, respectively. A focus error signal FES is given by the difference between the outputs of the second optical detector D2 and the third optical detector D3:FES=S2−S3 A tracking error is detected by a so-called three-beam method. The tracking subbeams are condensed onto the optical detectors D1 and D5. A tracking error signal TES is given by the difference between the outputs of the optical detectors D1 and D5:TES=S1−S5 A reproduction signal RF is given by the sum of the outputs of the second, third and fourth optical detectors D2,D3 and D4:RF=S2+S3+S4 In a hologram optical pickup using the conventional three-beam method, the hologram 4 includes two regions 4a and 4b having different grating pitches. Light beams which pass through the regions 4a and 4b of the hologram 4 after reflection on the optical disk 1 have different diffraction angles. Therefore, the light beams which have passed through the regions 4a and 4b are diffracted at a smaller angle and a larger angle in one direction with respect to the hologram 4. The grating of the hologram 4 is typically formed of grooves which are formed by patterning using a photoetching technique. When the two regions 4a and 4b having different grating pitches are formed by patterning, the depth of the grooves and the angle of grating vary in each of the regions 4a and 4b, depending on an etching rate, etc. If the groove depth and grating angle vary in each of the regions 4a and 4b, the variations appear as the difference in the intensity of diffracted light between the main beam and the subbeams, i.e., the difference in diffraction efficiency. As a result, the optical intensities of reflected light beams entering the optical detectors D1 to D5 are unbalanced, so that offset develops in the tracking error signal TES. In this case, characteristics of the hologram optical pickup are likely to be degraded.
{ "pile_set_name": "USPTO Backgrounds" }
The present invention relates to a new and distinctive soybean variety, designated RJS38003 which has been the result of years of careful breeding and selection as part of a soybean breeding program. There are numerous steps in the development of any novel, desirable plant germplasm. Plant breeding begins with the analysis and definition of problems and weaknesses of the current germplasm, the establishment of program goals, and the definition of specific breeding objectives. The next step is selection of germplasm that possess the traits to meet the program goals. The goal is to combine in a single variety an improved combination of desirable traits from the parental germplasm. These important traits may include but are not limited to higher seed yield, resistance to diseases and insects, tolerance to drought and heat, altered fatty acid profile, abiotic stress tolerance, improvements in compositional traits and better agronomic qualities. These processes, which lead to the final step of marketing and distribution, can take from six to twelve years from the time the first cross is made. Therefore, development of new varieties is a time-consuming process that requires precise forward planning, efficient use of resources, and a minimum of changes in direction. Soybean (Glycine max), is an important and valuable field crop. Thus, a continuing goal of soybean breeders is to develop stable, high yielding soybean varieties that are agronomically sound. The reasons for this goal are to maximize the amount of grain produced on the land used and to supply food for both animals and humans. To accomplish this goal, the soybean breeder must select and develop soybean plants that have the traits that result in superior varieties. Pioneer soybean research staff create over 500,000 potential new varieties each year. Of those new varieties, less than 50 and more commonly less than 25 are actually selected for commercial use. The soybean is the world's leading source of vegetable oil and protein meal. The oil extracted from soybeans is used for cooking oil, margarine, and salad dressings. Soybean oil is composed of saturated, monounsaturated and polyunsaturated fatty acids. It has a typical composition of 11% palmitic, 4% stearic, 25% oleic, 50% linoleic and 9% linolenic fatty acid content (“Economic Implications of Modified Soybean Traits Summary Report”, Iowa Soybean Promotion Board & American Soybean Association Special Report 92S, May 1990). Changes in fatty acid composition for improved oxidative stability and nutrition are also important traits. Industrial uses for processed soybean oil include ingredients for paints, plastics, fibers, detergents, cosmetics, and lubricants. Soybean oil may be split, inter-esterified, sulfurized, epoxidized, polymerized, ethoxylated, or cleaved. Designing and producing soybean oil derivatives with improved functionality, oliochemistry, is a rapidly growing field. The typical mixture of triglycerides is usually split and separated into pure fatty acids, which are then combined with petroleum-derived alcohols or acids, nitrogen, sulfonates, chlorine, or with fatty alcohols derived from fats and oils. Soybean is also used as a food source for both animals and humans. Soybean is widely used as a source of protein for animal feeds for poultry, swine and cattle. During processing of whole soybeans, the fibrous hull is removed and the oil is extracted. The remaining soybean meal is a combination of carbohydrates and approximately 50% protein. For human consumption soybean meal is made into soybean flour which is processed to protein concentrates used for meat extenders or specialty pet foods. Production of edible protein ingredients from soybean offers a healthy, less expensive replacement for animal protein in meats as well as dairy-type products.
{ "pile_set_name": "USPTO Backgrounds" }
The catalytic epoxidation of olefins using a silver-based catalyst has been known for a long time. Conventional silver-based catalysts have provided the olefin oxides notoriously in a low selectivity. For example, when using conventional catalysts in the epoxidation of ethylene, the selectivity towards ethylene oxide, expressed as a fraction of the ethylene converted, does not reach values above the 6/7 or 85.7 mole-% limit. Therefore, this limit has long been considered to be the theoretically maximal selectivity of this reaction, based on the stoichiometry of the reaction equation7C2H4+6O2=>6C2H4O+2CO2+2H2O,cf. Kirk-Othmer's Encyclopedia of Chemical Technology, 3rd ed., Vol. 9, 1980, p. 445. Modern silver-based catalysts however are highly selective towards olefin oxide production. When using the modern catalysts in the epoxidation of ethylene the selectivity towards ethylene oxide can reach values above the 6/7 or 85.7 mole-% limit referred to, for example 88 mole-%, or 89 mole-%, or above. Such highly selective catalysts, which may comprise as their active components silver, rhenium, at least one further metal and optionally a rhenium co-promoter, are disclosed in U.S. Pat. No. 4,761,394, U.S. Pat. No. 4,766,105, EP-A-266015 and in several subsequent patent publications. The silver based catalysts are subject to an aging-related performance decline during normal operation and they need to be exchanged periodically. The aging manifests itself by a reduction in the activity of the catalyst. Usually, when a reduction in activity of the catalyst is manifest, the reaction temperature is increased in order to compensate for the reduction in activity. The reaction temperature may be increased until it becomes undesirably high, at which point in time the catalyst is deemed to be at the end of its lifetime and would need to be exchanged. A reaction modifier, for example an organic halide, may be added to the feed to an epoxidation reactor for increasing the selectivity (cf. for example EP-A-352850). The reaction modifier suppresses the undesirable oxidation of olefin or olefin oxide to carbon dioxide and water, relative to the desired formation of olefin oxide, by a so-far unexplained mechanism. The optimal quantity of the reaction modifier depends on the epoxidation reaction conditions and on the type of catalyst used. Conventional catalysts have relatively flat selectivity curves for the modifier, i.e. the curves of the selectivity as a function of the quantity of the reaction modifier show that the selectivities are almost invariant over a wide range of reaction modifier quantities, and this property does virtually not change as a function of the reaction temperature and during prolonged operation of the catalyst. Therefore, when using a conventional catalyst, for optimum selectivity the quantity of the reaction modifier can be chosen rather freely and it can remain substantially the same during the entire lifetime of the catalyst. By contrast, the highly selective catalysts tend to exhibit relatively steep selectivity curves for the modifier, viz. for the highly selective catalysts the selectivity varies considerably with relatively small changes in the quantity of the reaction modifier, and the selectivity exhibits a pronounced maximum, i.e. an optimum, at a certain quantity of the reaction modifier. This has been illustrated in EP-A-352850 (cf. FIG. 3 therein). Moreover, the selectivity curves and more in particular this quantity of the reaction modifier where the selectivity is at optimum tend to change with the reaction temperature and, thus, during the catalyst life. Consequently, when employing the highly selective catalysts in combination with a reaction modifier, the selectivity may vary to an undesirably large extent with changes of the reaction temperature and over the lifetime of the catalyst. Namely, when the reaction temperature is changed, for example to compensate for a reduction in the activity of the catalyst, it represents itself as a problem to maintain reaction conditions which are optimal with respect to the selectivity towards the olefin oxide production.
{ "pile_set_name": "USPTO Backgrounds" }
The present invention relates to an apparatus, namely a sheet, comforter and pillow cover, which is designed to provide a sanitary sleeping environment away from home. The sheet cover will also necessitate comfort and ease of mind by ensuring that the remnants of disease, bacteria, viruses and body contaminates which may be present from past occupants, will not come in contact with the individual utilizing the enclosed apparatus.
{ "pile_set_name": "USPTO Backgrounds" }
Reference to any prior art in the specification is not, and should not be taken as, an acknowledgment or any form of suggestion that this prior art forms part of the common general knowledge in Australia or any other jurisdiction or that this prior art could reasonably be expected to be ascertained, understood and regarded as relevant by a person skilled in the art. The separation of CO2 from flue gases in power stations, cement kilns and in steel manufacturing allows these industrial activities to proceed with the use of fossil fuels, while reducing the emissions of the most important greenhouse gas, i.e. CO2. Although several different processes are currently under development for the separation of CO2 from flue gases, chemical absorption processes using aqueous solutions of chemical absorbents are the leading technology, mainly because of its advanced state of development. While it is already available at low CO2-removal capacities, it is not at the scale necessary for large scale industrial operation. Scaling up the process is therefore a major challenge. The typical flow sheet of CO2 recovery using chemical absorbents is shown in FIG. 1. After cooling, the flue gas 100 is brought into contact with the chemical absorbent in an absorber 102. A blower 103 is required to pump the gas through the absorber after passing through a cooler 104 at temperatures typically between 40 and 60° C. whereby CO2 is then bound by the chemical absorbent in the absorber 102. After passing through the absorber 102, the flue gas undergoes a water wash section 113 to balance water in the system and to remove any droplets of vapour carried over and then leaves the absorber 102. The “rich” absorbent solution, which contains the chemically bound CO2, is then pumped to the top of a stripper 107, via a heat exchanger 105. The regeneration of the chemical absorbent is carried out in the stripper 107 at elevated temperatures (100-140° C.) and pressures between 1 and 2 bar. The stripper 107 is a gas/liquid contactor in which the rich absorbent is contacted with steam produced in a reboiler 108. Heat is supplied to the reboiler 108 to maintain the regeneration conditions. This leads to an energy penalty as a result of the heating up the solution to provide the required desorption heat for desorbing the chemically bound CO2 and for steam production which acts as a stripping gas. Steam is recovered in a condenser 109 and fed back to the stripper 107, whereas the CO2 product gas 110 leaves the condenser 109. The heat of condensation is carried away in cooling water or an air cooling device. The CO2-product 110 is a relatively pure (>99%) product, with water vapour being the main other component. Due to the selective nature of the chemical absorption process, the concentration of inert gases is low. The “lean” absorbent solution 111, containing far less CO2 is then pumped back to the absorber 102 via the lean-rich heat exchanger 105 and a cooler 112 to bring it down to the absorber temperature level. CO2 removal is typically around 90%. The energy requirement of a chemical absorption process mainly stems from the heat supplied to the reboiler 108. This heat is used to produce steam from the lean solution 106 which acts a stripping gas, i.e. it keeps the partial pressure of CO2 sufficiently low to provide a driving force for the desorption process. The steam is also the carrier of thermal energy which, through its condensation, releases the energy required to desorb CO2 and to heat up the chemical absorbent through the desorption column 107. The amount of steam generated in the reboiler 108 should be kept as a low as possible, but some of the steam will always inevitably be lost from the desorption unit with the CO2-produced and this represents an energy loss, as the steam is usually condensed and the energy is carried away in the cooling water. There exists a need to provide a process and apparatus that is more energy efficient than the present process. Various approaches have been suggested in the prior art. WO 2007/075466 discloses an approach in which the lean chemical absorbent exiting from the bottom of the desorption column, is sent to a flash vessel at lower pressure than the equilibrium pressure. The lean solution has a low CO2 partial pressure; hence the vapour produced is predominantly steam. The steam can be recompressed and injected into the desorber to provide additional steam for stripping and heating of the chemical absorbent in the desorption column. This requires the addition of two pieces of equipment: flash vessel and a compressor. This adds to capital and operating expenses. Furthermore, the addition of a further compressor adds to the overall energy consumption—thus detracting from energy savings, which is undesired. U.S. Pat. No. 4,152,217 discloses the recovery and reuse of heat by using a heat exchanger to extract heat energy from the CO2/steam mixture exiting the desorption column through heat exchange with the rich absorbent solution entering the desorption column. This allows the latent heat of condensation to be recovered as sensible heat in the rich absorbent. The teachings of this patent show however that a split of the rich absorbent flow into two streams is needed, which is complicated from a process flow and control perspective. U.S. Pat. No. 4,444,571 discloses a method to recover latent heat from gas/steam mixtures that uses a membrane which is selectively permeable towards steam over other gases contained in the mixture. Permeate which predominantly contains steam is recompressed and injected into the bottom section of the desorption column. Although suitable membranes are available, the process requires additional energy, which is undesirable. The state-of-the-art methods to improve the energy performance of the desorption column are all limited because they only address a single improvement step and assume the usual process lay-out as shown in FIG. 1. It is an object of the invention to address at least some of the above aforementioned short-comings of the prior art.
{ "pile_set_name": "USPTO Backgrounds" }
Reconstructing 3D scenes from 2D images is a very important problem in computer vision and other imaging applications. Conventional 3D reconstruction methods typically use two or more images to obtain depths in the scene. However, depth recovery from a single 2D imago is a severely ill-posed problem. Rather than reconstructing a 3D scene using geometric entities, such as points and polygons, one method uses a 3D reconstruction procedure that constructs a popup model. Using several image and geometric features, that method automatically classifies regions as ground, buildings and sky. Another method infers absolute depth using image features and weak assumptions based on coplanarity and connectivity constraints. For modeling indoor scenes, one method uses a cuboid model to approximate geometry of a room. With that model, pixels in an image are classified as left wall, middle wall, right wall, floor and ceiling. For indoor scenes, we refer to this classification as the indoor scene layout estimation or just layout estimation. To estimate the optimal layout, hundreds of cuboids are sampled and each cuboid is assigned a score based on several image and geometric features. We refer to this cuboid estimation problem as layout estimation. That method uses training images to classify texture, color and line features to obtain the pixel-level classification.
{ "pile_set_name": "USPTO Backgrounds" }
The present invention relates to a nuclear magnetic resonance spectrometer suitable for analyzing, in a liquid-solution, the structure and interaction of protein and organic molecules such as substrate and ligand interacting with the protein. A method for analyzing organic matter utilizing nuclear magnetic resonance (NMR) has been making rapid progress in recent years. In particular, a method has been used in combination with the technique of strong superconductive magnets to make it possible to highly efficiently analyze the structure of an organic compound such as protein having a complicated molecular structure on atomic level. The present invention is directed to a nuclear magnetic resonance (NMR) spectrometer necessary for analyzing the structure on atomic level and interaction of protein molecules in an aqueous solution dissolving a small quantity of protein, especially, an energy spectrometer differing from a medical MRI diagnostic apparatus aiming at imaging of a tomogram of human body needing so-called millimeter class image resolution in that performance such as the magnetic field intensity being one order or more higher, the magnetic field uniformity being of four order and the stability being three order higher is needed, thus requiring quite different design technique and apparatus manufacture technique. A detailed description of a conventional high-resolution NMR spectrometer is given in “NMR for Protein” by Yoji Arata, published by Kyoritu-shuppan, pp. 33–54, 1996. As up-to-date inventions concerning the typical apparatus construction when the NMR is utilized for analysis of protein, one may refer to an invention relating to a superconductive magnet that shows the typical construction of a multilayer air-core solenoid coil as disclosed in JP-A-2000-147082, an invention relating to a signal detection technique that shows a cage type superconductive detection coil as disclosed in U.S. Pat. No. 6,121,776 and examples of signal detection technique based on a conventional barrel type or cage type coil as disclosed in JP-A-2000-266830 and JP-A-6-237912. According to these reports, all of the conventional high-sensitivity NMR spectrometers for protein analysis use a superconductive magnet unit constructed of solenoid coils used in combination to generate a magnetic field in the vertical directions, an electromagnetic wave at 400 to 900 MHz is irradiated on a sample and a resonance wave generated from the sample is detected by utilizing the barrel or cage type detection coil. Also, as shown in an example of U.S. Pat. No. 6,121,776, a contrivance is made to improve the S/N sensitivity ratio by utilizing a detector cooled to low temperatures with a view to decreasing heat noise during reception. Historically, in the high-sensitivity NMR spectrometer, improvements in sensitivity have been achieved by a method in which the basic constituents of a system such as antenna and magnet are kept to remain unchanged but the center magnetic field intensity of the superconductive magnet is increased. Accordingly, the maximum NMR measurement sensitivity reported till now can be obtained with a NMR spectrometer at 900 MHz utilizing a large superconductive magnet having a center magnetic field of 21.1 tesla and the basic constituents of the spectometer remain unchanged as compared to those in the prior art of JP-A-2000-147082. In the analysis of protein using a liquid-solution, the improved center magnetic field is effective to clarify separation of improvement in sensitivity and chemical shift. In attaining the effect of improved sensitivity attributable to the form or shape of detection coil, it has been known, for example, as described in “Book of NMR” by Yoji Arata, published by Maruzen, PP. 325–327, 2000, that the solenoid coil conventionally used as detection coil is advantageous over the barrel or cage type in various points. For example, the solenoid coil is advantageous in easy controllability of impedance, filling factor and efficiency of RF magnetic field. But, according to the literature, when the sensitivity is thought much of in such an application of measuring protein dissolved by a small quantity in an aqueous solution, winding a solenoid coil around a sample tube placed vertically to the magnetic field is practically impossible and in general, is not utilized. In an exceptional application where highly sensitive measurement is carried out by using a small quantity of sample solution, the above technique is limitedly utilized through a method utilizing a particularly designed micro-sample tube to carry out measurement by using a special probe. Recently, in a special example disclosed in JP-A-11-248810, a method is contrived according to which a bulky magnet of high-temperature superconductivity is magnetized in the horizontal direction and a NMR signal is detected with a solenoid coil. Further, JP-A-7-240310 discloses a method for constructing a superconductive magnet and a cooling container to meet a general NMR application directed to eliminate constraints on the top or ceiling height of apparatus. However, any method of improving the detection sensitivity necessary for analyzing protein and any method of technically coping with magnetic field uniformity and temporal stability of magnetic field have not been known yet. Recently, with needs for study of protein promoted, needs for analyzing a sample in which protein has a small degree of solubility to water have increased and there is a need of improving the sensitivity of measurement of NMR. To adapt the NMR spectrometer to the needs as above, the measurement sensitivity must be improved while maintaining a sample space comparable to that of the conventional apparatus and besides the maintenance of the stability of a superconductive magnetic field over a long time of data integration is indispensable. The improved measurement sensitivity is particularly advantageous in that for samples having substantially the same solubility, not only the measurement time can be shortened but also the sampling amount can be decreased, thereby ensuring that the protein of small solubility can be analyzed to advantage. Accordingly, the NMR spectrometer used for analysis of protein is required of far more excellent detection sensitivity and stability than those in the conventional NMR and in addition, is required to have ability to detect NMR signals accurately and stably over a long time of one week or more. This is because if the magnetic field varies during measurement, the peak of NMR signal is caused to shift and especially, in measurement of interaction, the peak shift due to interaction cannot be discriminated from that due to instability of the magnetic field. If the magnetic field is non-uniform, desired peaks overlap each other, raising a problem that discrimination of interaction is difficult to achieve. Therefore, it should be noticed that, in future NMR techniques aiming at performing various kindes of analysis of protein, development of new technology not lying on mere extension of the conventional general NMR spectrometers will be needed. For example, specification of magnetic field uniformity in the general NMR spectrometer is 0.01 ppm in a sample space, that is, 0.01 ppm in terms of temporal stability. When this value is indicated in terms of proton NMR for general 600 MHz use, a permissible error of 6 Hz results. In the case of the aforementioned analysis of interaction of protein, however, spatial and temporal resolution of at least 1.0 Hz or less is required and preferably, 0.5 Hz or less is needed. In a method capable of implementing the magnetic field stability and the temporal stability of magnetic field, the construction of superconductive magnet and detection coil must be optimized. Accordingly, the performance of the conventional, generally-used NMR spectrometer is insufficient and the stability and magnetic field uniformity higher by one order or more than those of the conventional spectrometer are required. In the prior arts, the sensitivity is managed to be improved by relying on improvements in magnetic field intensity and as a result, the apparatus is increased in size and to cope with problems of leakage magnetic field and floor strength, there arises a new problem of installation capability such as needs for a dedicated building. Further, disadvantageously, the cost of a superconductive magnet increases. The improved sensitivity has an upper limit of about 21 T because of constraints due to a critical magnetic field of a superconductive material and for more upgraded improvements in sensitivity, the advent of a technique for improving detection sensitivity based on a new means without resort to the magnetic field intensity has been desired. The aforementioned high-sensitivity measuring method utilizing the solenoid coil can be used with a special sample tube for a very small quantity of sample and a special detection probe but it cannot be applied to analysis based on a general protein solution of about 10 cc. The method for generating a magnetic field in the horizontal direction by means of a strong magnet and detecting NMR signals by means of a solenoid coil as described in the example of JP-A-11-248810 can generate only a magnetic field of not greater than 10 T at the surface of a high-temperature superconductor, with the result that the magnetic field at a sample part is about several tesla at the most, thus proving that the method of interest cannot generate a magnetic field of 11 tesla or more necessary for analysis of protein, preferably, a magnetic field of 14.1 tesla or more in a desired sample space. Further, in this method, owing to the effect of a magnetic flux creep phenomenon of the high-temperature superconductor, the temporal stability 1.0 Hz/hour or less necessary for analysis of protein is substantially difficult to achieve. As regards the magnetic uniformity necessary for analysis of protein, non-homogeneity attributable to the manufacture process of a high-temperature superconductive bulky material also makes it difficult to attain the magnetic field uniformity within 1.0 Hz in terms of proton NMR frequency in a space defined by 10 mm diameter×20 mm length. As described above, while a breakthrough technique meeting the needs for analysis of protein is desired to be developed in connection with the conventional techniques, the advent of a new solving method for further improvements in sensitivity has been desired under the present-day circumstances that improving the sensitivity based on the magnetic field has reached limits. For the purpose of conducting an efficient and accurate. Analysis of the interaction of protein in a liquid-solution with low molecules such as substrate and ligand, for which needs are considered to increase in future, it is empirically preferable that a suitable quantity of sample be measured at 600 to 900 MHz and with a center magnetic field of about 14 to 21 T and the measurement sensitivity be increased beyond the present one to increase the throughput. Generally, in a spectrometer operating at 800 MHz or more, for the purpose of making full use of the superconductive characteristics to an extreme, operation is carried out by depressurizing liquid helium at 4.2 K and excessively cooling it to 1.8 K. Therefore, complexities in apparatus operation are aggravated and maintenance is laborious. In addition, the magnetic unit increases in size to increase leakage magnetic field and typically, a dedicated building is needed. Especially, the leakage magnetic field in the vertical direction increases as the center magnetic field increases in the conventional system, so that in an apparatus of 900 MHz class, for instance, a leakage magnetic field occurs extending up to 5 m in the height direction, and from the viewpoint of apparatus installation, there needs a tall building of high ceiling. As a result, the construction cost increases disadvantageously. Further, the conventional 900 MHz superconductive magnet is sized such that only a magnet part has a diameter of 1.86 m and a height of several meters, as described in IEEE. Transactions on Applied Superconductivity, Vol. 10, No. 1, page 728–731. The present invention intends to provide a novel NMR spectrometer in which the measurement sensitivity of NMR signals can be increased by at least 2.5 times or more of that in the conventional apparatus at about 600 MHz (14.1 T) under a condition that a normal sample tube of 5 to 10 diameter is mainly used and a sample liquid-solution is charged in the tube up to a height of about 30 mm and the temporal stability and spatial uniformity of a superconductive magnet necessary for analysis of protein can be provided. In the construction of the present invention, the operating temperature of a system is not set to 4.2 K. By applying the present invention, it is also possible to aim at achieving extremity performance but depending on applications, operation at the conventional magnetic field limit 21.1 T, that is, at 900 MHz and at 1.8 K can proceed and in that case, the sensitivity can be improved by 40% of that in the conventional system, proving that overcoming the detection sensitivity limit attributable to magnetic field intensity, conventionally unattainable, can succeed for the first time.
{ "pile_set_name": "USPTO Backgrounds" }
The present invention relates to the detection of explosives hidden in packages, particularly small amounts of modern, highly-explosive, nitrogen-based plastic explosives hidden in airline bags. The detection of explosive devices hidden in airline baggage is a significant problem, particularly in view of the development of modern plastic explosives which can be formed into various innocent-appearing shapes and which are sufficiently powerful that small quantities can destroy an aircraft in flight. In general, three different screening approaches for the detection of hidden explosives are known and employed to various degrees at certain airports. The first of these approaches is conventional X-ray imaging. Mere X-ray imaging however is of limited effectiveness, particularly since explosives need not be formed into any particular shape. The second approach is the use of a so-called vapor sniffer which collects vapors emanating from luggage and analyzes them for the presence of molecules of explosive materials. While such devices are relatively sensitive, they nevertheless cannot detect explosives which are sealed within containers so as to prevent the escape of sufficient vapor quantities for detection. The third approach involves the detection of nitrogen by means of thermal neutron interrogation. Nitrogen is a component of virtually every practical known high explosive. Thermal neutron interrogation involves exposing baggage to a "sea" of thermal neutrons (or "slow" neutrons having an energy in the order of 0.025 eV). Thermal neutrons combine with the nuclei of nitrogen-14 atoms to produce an energetic form of nitrogen-15 isotope. The energetic nitrogen-15 isotope immediately decays to its ground state, emitting characteristic 10.8 MeV gamma rays in the process. The 10.8 MeV gamma rays are detected as indicator of the presence of nitrogen in the package. There are, however, a number of problems with such detectors employing thermal neutrons. A typical neutron source is radioactive californium-252 which emits energetic neutrons that are then slowed to thermal energies for reaction with nitrogen-14 nuclei. The use of such a radioactive neutron source introduces logistical problems related to handling and radiation shielding. Other disadvantages of such systems include a relatively high false-positive rate coupled with the inability to effectively detect small quantities of explosives, such as quantities less than about one or two pounds. In particular, there are a number of materials other than explosives which contain nitrogen, such as wool and leather. If the threshold level of a system employing thermal neutron interrogation is adjusted so as to detect small quantities of nitrogen, then a high false-positive rate results due to the presence of innocent nitrogen-containing materials, leading to the necessity of searching an excessive number of packages by hand, negating the practical effectiveness of the system. If, on the other hand, the threshold level is set high to avoid false-positives, then the likelihood that actual explosives will escape detection is increased, again negating the effectiveness of the system.
{ "pile_set_name": "USPTO Backgrounds" }
Embodiments described generally relate to the field of communication systems and, more particularly, to enhanced wireless network coverage within a defined coverage area. Wireless networks are widely deployed, especially for use in well-defined and limited areas such as homes and apartments. Often, a single wireless access point/router is connected to a cable modem or digital subscriber line (DSL) modem to provide wireless access to a broadband network. Wireless access points can operate in a 2.4 GHz frequency band or in both the 2.4 GHz frequency band and a 5 GHz frequency band. However, these wireless access points may not be able to support the demands of multiple end devices (or stations). In particular, the limited frequencies used by the wireless access points may not support high data throughput rates required by multiple stations.
{ "pile_set_name": "USPTO Backgrounds" }
1. Field of the Invention The present invention relates to a zoom lens usable as a photographing lens of a lens shutter camera and a photographing lens of a video camera. 2. Description of the Related Art Recently, a zoom lens has been mostly mounted onto a lens shutter camera. A new problem of how to make the camera having the zoom lens compact is caused. When the camera is made compact, a maximum problem from a viewpoint of portability of the camera is to reduce a size of the camera in an optical axis direction of the zoom lens. Namely, this maximum problem is to make the zoom lens thin. Recently, the lens shutter camera mounting the zoom lens thereon is made thin by increasing the number of stages in body tube. However, it is important to make the zoom lens itself thin. In particular, an entire length of the lens shutter camera at a long focal end (a distance from a lens front end to an image face) and an entire thickness of all movable lens groups at a most approaching time thereof (a distance from a lens front end to a lens rear end) are most important as a factor for determining thinness of the camera. A zoom lens constructed by two lens groups as shown in Japanese Patent Application Laying Open (KOKAI) No. 5-11186 is generally known as a zoom lens having a relatively short length as the above entire lens length. A zoom lens constructed by two lens groups as shown in Japanese Patent Application Laying Open (KOKAI) No. 5-127082 is known as a zoom lens having a relatively thin thickness as the entire lens group thickness. An entire length of the former zoom lens ranges from 0.96 to 0.92 times the focal length of an entire lens system at a long focal end so that this entire length of the former zoom lens is short. However, the entire lens group thickness of the former zoom lens is about 1.25 times the height of a maximum image. This entire lens group thickness is thick and is still insufficient to make the camera thin. An entire lens group thickness of the latter zoom lens at a long focal end is approximately equal to the height of a maximum image so that the latter zoom lens is suitable for making the camera thin. However, a zoom ratio as an important factor is a small value such as about 1.6 magnifications so that the latter zoom lens is not necessarily practical.
{ "pile_set_name": "USPTO Backgrounds" }
When operating an on-vehicle disk brake lathe, it is beneficial to be able to monitor the cutting status to establish the likelihood of a continuous cut being performed. To track the cutting status of a brake lathe during the turning operation, the lathe should be provided with a system for making a determination of when tool bits of the lathe have begun actively cutting surfaces of a brake disk to machine them, as well as a determination of when such cutting, once initiated, has ceased. Past systems for monitoring the contact of tool bits of a lathe with the surfaces of a brake disk during a turning operation have employed contact sensors to determine when the tool bits of the lathe are in contact with the disk surfaces, in order to facilitate setting an appropriate depth of cut, such as taught in U.S. Pat. No. 6,363,821. This patent teaches that such contact sensors can include electrical continuity detectors, vibration sensors, and strain sensors. Electrical conductivity sensors and vibration sensors are also taught in U.S. Publication 2005/0016338 for determining when the tool bits of a lathe are in contact with disk surfaces. The use of an electrical conductivity sensor to determine when a cutting operation has been completed is taught in U.S. Pat. No. 6,729,212, which teaches an iterative machining technique where machining is repeated with increasing depth of cut until such time as the electrical contact between the tool bit and the disk surface is determined to be substantially constant throughout the cut.
{ "pile_set_name": "USPTO Backgrounds" }
This invention relates to an exhaust control valve system for a parallel multi-cylinder, two-cycle engine and more particularly to an improved exhaust control valve system for such an engine that permits a compact engine construction and near optimum running under all engine conditions. As is noted in our copending application of the same title, Serial No. 07/566,968, filed Aug. 13, 1990, and assigned to the assignee hereof, the performance of an internal combustion engine operating on the two-cycle, crankcase compression principle can be improved through the use of an exhaust control valve that changes the port configuration and timing in response to engine running characteristics. That is, it is desirable to provide a large overlap to achieve high speed performance but such large overlaps reduce the performance at low and mid ranges. The use of an exhaust control valve can change the timing of the exhaust port so as to optimum at all running conditions. In connection with the application of this principle, however and as noted in our aforenoted copending application, the positioning of the scavenge and exhaust ports can give rise to greater than desired engine length with inline engines. As also noted in that application, this problem can be overcome by rotating the ports about the cylinder bore axis so as to permit the ports of adjacent cylinders to nest between each other and thus reduce the engine length. However, this means that the exhaust port will then not extend in a direction perpendicularly to a plane containing the cylinder bore axes but is disposed at an acute angle to such a perpendicular plane. In the copending application, the control valves for the exhaust ports are mounted on shafts that extend perpendicularly to the center line of the exhaust passages and thus will provide a symmetrical port configuration regardless of the degree of opening or closure of the exhaust control valve. The construction as employed in our aforenoted copending application, however, presents certain difficulties in valve actuation. Since the shafts are all parallel to each other rather than aligned, a somewhat complicated arrangement must be employed so as to insure that all valves will be positioned similarly in response to the varying engine condition. Although it is possible to achieve this result, the structure for doing so tends to become complicated. It is, therefore, a principal object of this invention to provide an improved and simplified and yet compact control valve arrangement for the exhaust ports of a multiple cylinder, two-cycle internal combustion engine. It is possible to provide an exhaust port configuration as described in our aforenoted copending application and wherein all of the control valves are affixed to a common shaft which does not extend perpendicularly to the axis of the exhaust ports. However, when this is done, the movement of the control valve element will cause a variation in the configuration or symmetry of the exhaust port depending upon the position of the control valve member. This can give rise to less than optimum performance. It is, therefore, a further object of this invention to provide a control valve arrangement for the exhaust port of a two cycle internal combustion engine of the type wherein the axis of rotation of the control valve does not extend perpendicularly to the axis of the exhaust port and wherein the exhaust port and control valve are configured so as to minimize the amount asymmetry that occurs during the pivotal movement of the control valve. It is, therefore, a still further object of this invention to provide an improved control valve arrangement for a multiple cylinder inline, two-cycle, crankcase compression engine having angular disposed exhaust passages and control valves that are operated by a common shaft that extends in a non-perpendicular relationship to the exhaust passages and wherein in at least one operative position of the control valve the exhaust port is symmetrical.
{ "pile_set_name": "USPTO Backgrounds" }
1. Field of the Invention The present invention relates to an imaging apparatus and in particular to a technique for correcting a displacement between images. 2. Description of the Related Art Various techniques have conventionally been provided for correcting a displacement between images when synthesizing a plurality of static images photographed by a digital camera. A technique for correcting a displacement between images due to an unsteady camera is suitable to processing an image photographed in the dark because an exposure time required for photographing therein is longer, increasing the occurrence of an unsteady camera. Specifically, there are methods of correcting an unsteady camera (also noted as “image stabilization method” hereinafter) by means of hardware and software. The image stabilization method by means of hardware includes a lens shift system which moves lens in the direction opposite to an actual moving direction of the camera, and a CCD (charge coupled device) shift system. As to a technique for stabilizing image by software, provided is a technique which reduces an input image in a plurality of steps, utilizes a global shift amount roughly obtained by using a reduced image on a lower layer level and builds up a scan block when obtaining a local shift amount on a higher level, thereby making it possible to minimize a scan block set for a second image (refer to a patent document 1 for example). Another provision is a technique capable of speeding up a process for synthesizing a plurality of digital image data by making and storing a thumbnail for each of synthetic images in the case of performing a synthesizing process by using a plurality of image data (refer to a patent document 2 for example). At this point, a description is a detail on a conventional calculation method of a displacement between images among image stabilization techniques by means of software. FIG. 1 is a flow chart showing a displacement amount calculation process according to a conventional technique. The conventional technique obtains, in memory, a plurality of image data acquired from an imaging apparatus and detects a feature point (i.e., an edge) in each image, thereby calculating the displacement amount of each image when making a synthetic image by calculating respective displacement amounts related to a plurality of images and correcting the displacements as shown in the flow chart of FIG. 1. Additionally, as to an image stabilization technique using software, provided is a technique of detecting an unsteady camera by means of software (refer to a patent document 3 for example). [Patent document 1] Laid-Open Japanese Patent Application Publication No. 2004-343483 [Patent document 2] Laid-Open Japanese Patent Application Publication No. 2004-234624 [Patent document 3] Laid-Open Japanese Patent Application Publication No. 2005-197911 The above noted lens shift system and CCD shift system are both faced with a problem of an increased cost because hardware must be added to a camera for moving hardware such as a lens, CCD, et cetera. Another problem associated with these methods is that means of the image stabilization is prone to a shock and easy to fail. As for an image stabilization method using software, individual images must be positioned with each other before compounding the first and second images as described above. Accordingly exercised is to detect a feature point from the first image and obtain a position corresponding to the feature point in the second image and calculate a displacement amount from the difference of positions between the images. Such a calculation method of a displacement amount requires the process regarding the entirety of the image for detecting a feature point and searching a position corresponding thereto. However, a memory capacity for storing a plurality of images tends to get larger in proportionate with digital cameras with bigger pixels, and it causing an increase of cost. Also associated with digital cameras with bigger pixels, processes for an image stabilizer increase, hence it creates a problem of increasing the processing time.
{ "pile_set_name": "USPTO Backgrounds" }
1. Field of the Invention The field of the invention relates to methods of feature extraction, such as edge detection. It can be used in computer vision systems, including image/facial/object detection/recognition systems, scene interpretation, classification and captioning systems. 2. Technical Background Most of the existing object detection algorithms are based on machine learning classifiers, which in their turn use features extracted from an image. Fundamentally there are two approaches to enhance the results of an object detection algorithm. The first approach is an enhancement of classification methodology, where many techniques have been proposed in the literature (Linear classifiers, Neural networks etc.). The second approach is an enhancement of features used. Researches who focus their work on the enhancement of features extracted from an image mostly concentrate on finding the set of discrete primitives describing the image content. The process of feature extraction is usually related to filtering of the image data and normalisation of the filter's response. However, there is one common flaw in most feature extraction techniques, i.e., during the normalisation and accumulation of image features the assumption is made that the filters producing a stronger response represents stronger image features. In practice, research is often being carried out with digital video or photographic images that are products of image processing pipelines, processing the image sensor data with unknown settings. As previously discussed, such processing can significantly alter image data, breaking linear dependencies between parts of an image and unbalancing the appearance of different image elements. This invention provides a solution for a more robust edge detection method by taking sensor characteristics into account during edge detection. The method may be used for feature extraction or feature detection. 3. Discussion of Related Art The research being conducted at present in the object detection and classification area is very intense. There are a number of object detection techniques, among which HOG-SVM and CNN are widely used. One of the most successful object detection techniques is known as Histogram of Oriented Gradients—Support Vector Machine (HOG-SVM) as described in [1-5]. The results produced by object detection algorithms are continuously improving. The first step in the calculation of Histogram of Oriented Gradients is edge detection. Standard approaches are presented in [6-10]. A Convolutional Neural Network (CNN) is a type of feed-forward artificial neural network (ANN) where the individual neurons are tiled in such a way that they respond to overlapping regions in the visual field. When used for image recognition, convolutional neural networks (CNNs) consist of multiple layers of small neuron collections which look at small portions of the input image, called receptive fields. The results of these collections are then tiled so that they overlap to obtain a better representation of the original image; this is repeated for every such layer. The layers form a hierarchical system in which the first layers look for lower level features; this is accomplished by means of convolution between a filter and an image. Existing approaches that assume the object detection algorithm will run on recorded video or still image present a number of issues. First, object detection always needs image processing system to produce a quality RGB image or video sequence, which in many cases means increased system complexity. Secondly the object detection algorithms assume no knowledge about image source, as the image processing settings are not known. Therefore the performance of object detection algorithms may deteriorate quickly in low light conditions.
{ "pile_set_name": "USPTO Backgrounds" }
Crowded areas, such as transit systems, malls, event centers, and the like often have security staff monitoring video systems to ensure that security and/or emergency situations are promptly detected and dealt with. However, in many cases, the crowded areas are too crowded or so large that it is difficult or impossible to monitor all of the necessary video feeds for a particular space without using an inordinately large staff of security personnel. Additionally, the human element involved in recognizing particular behaviors can result in delays in recognizing and responding to such behavior. Improvements in such security and emergency scenario recognition are desired.
{ "pile_set_name": "USPTO Backgrounds" }
1. Field of the Invention The present invention is directed to the annotation of source code to facilitate source code comprehension, particularly for software development. More particularly, the invention concerns the use of annotation information for source code debugging. 2. Description of the Prior Art By way of background, conventional software debugging tools allow a user to monitor program execution, check interrupt and exception handling, inspect data values, and perform other tasks that allow software operation to be checked and verified. Using such tools requires intimate knowledge of the source code, including the executable file structure and the function and method names for the program. This is because debugging typically involves halting program execution by setting breakpoints at selected instructions or by setting watchpoints that trigger breakpoints when specified memory addresses are accessed or reach certain values or when subroutines or other sections of code are executed. Thus, when a developer is required to debug unfamiliar source code, he/she is faced with the need to understand what may be a large and complex code base. Existing annotation tools are available that describe functions, methods, or code blocks using source code comments, state diagrams, flow charts, call graphs, control flow graphs, pseudo code, and/or a description of what a specified code block accomplishes. However, even if extensive annotation is available, the developer must ultimately wade into the code in order to perform debugging tasks such as setting break points, watch points, etc. so that the desired portion(s) of the software may be observed during execution. It is to improvements in the art of software debugging that the present invention is directed. In particular, a technique is described wherein software metadata resources can be leveraged to directly support the debugging effort.
{ "pile_set_name": "USPTO Backgrounds" }
The subject matter discussed in the background section should not be assumed to be prior art merely as a result of its mention in the background section. Similarly, a problem mentioned in the background section or associated with the subject matter of the background section should not be assumed to have been previously recognized in the prior art. The subject matter in the background section merely represents different approaches, which in and of themselves may also be inventions. An on-demand system typically includes multiple nodes, such as application servers, that interact with a relational database management system. Each of the nodes may convey selected data updates from a relational database to multiple clients that subscribe to the on-demand system. Each node typically accesses the relational database management system to identify each of the data updates based on timestamps corresponding to the data updates. For example, after an application server conveys selected data updates from the relational database to subscribing clients beginning at 9:00 AM, the relational database management system receives data updates at 9:19 AM, and the application server checks for additional data updates at 9:30 AM by examining the timestamps of the data updates in a transaction table stored in the relational database. The application server skips the data updates with a timestamp from before 9:00 AM because the application server presumably processed these data updates already beginning at 9:00 AM, and only attempts to process the data updates with timestamps between 9:00 AM and 9:30 AM. However, the relational database management system may receive data updates in bursts, such as 2 million data updates between 9:00 AM and 9:30 AM at 9:19 AM, and no data updates between 9:30 AM and 10:00 AM. Although the relational database management system may have the capacity to receive and process millions of data updates in a short period of time, an application server may not have the capacity to process millions of data updates before the application server is scheduled to process data updates again. Therefore, when the application server is scheduled to process data updates again at 10:00 AM, the application server may not have completed the processing of the 2 million data updates that the application server began processing at 9:30 AM. In this situation, the application server may malfunction, and even crash. The application server may terminate the 9:30 AM job before all of the 2 million data updates were processed, such that the application server does not process the unfinished data updates because these unfinished data updates have 9:19 AM timestamps and the 10:00 AM job only processes updates with timestamps after 9:30 AM, when all of the previous data updates were presumed to have already begun processing. Even though many data updates remain unprocessed, the application server may be idle because no data updates exist with a timestamp between 9:30 AM and 10:00 AM. Alternatively, the application server may begin reprocessing some of the data updates that the application server already processed beginning at 9:30 AM, which may not only waste resources, but also result in yet another job timeout before all of the data updates are processed. Further problems may exist when an application server is taken offline for scheduled maintenance or to repair a problem. When the application server is bought back online, the timestamp of the application server's most recent job may be so long ago relative to the recent data updates that the application server may experience many problems, including some of the problems described above, in an attempt to process all of the data updates that occurred during the application server's downtime.
{ "pile_set_name": "USPTO Backgrounds" }
Various types of swivel mechanisms have been known heretofore for use with standard and recliner outdoor chairs. These prior mechanisms are exemplified for example in U.S. Pat. Nos. 2,687,765, 2,914,793 or 637,968 of record in the copending parent application. Such prior types of swivel chairs and mechanisms however have been poorly adapted for ordinary outdoor use as, for example, on lawns, near pools or on sandy surfaces. Lightweight outdoor chairs with swivel mechanism have not been available heretofore which enable people easily to carry the chair to a convenient outdoor location, to use the chair, to return the chair to storage and to clean the swivel mechanism when necessary without undue effort. In addition, such prior swivel chairs and mechanisms have not been known to provide adequate stability, especially on uneven terrain. Moreover, swivel mechanisms have not been available heretofore which are easily mountable to and demountable from a standard lightweight outdoor chair or recliner thereby enabling the conversion of virtually any such standard outdoor chair or recliner into a swivel chair. Finally, outdoor swivel structures have not been heretofore available which will function properly even when constantly exposed to dirt or sand and which may be easily disassembled for suitable cleaning without undue effort. This has been due in part to the fact that outdoor furniture is usually relatively low cost furniture and suitably low cost swivel mechanisms have not heretofore been available for outdoor use.
{ "pile_set_name": "USPTO Backgrounds" }
Automotive and many industrial lubricating oils are usually formulated from paraffin based petroleum distillate oils or from synthetic base lubricating oils. Lubricating oils are combined with additives such as soaps, extreme pressure (E.P.) agents, viscosity index (V.I.) improvers, antifoamants, rust inhibitors, antiwear agents, antioxidants, and polymeric dispersants to produce an engine lubricating oil of SAE 5 to SAE 60 viscosity. After use, this oil is collected from truck and bus fleets, automobile service facilities, municipal motor oil recycling centers and retail stores. There is also a significant volume of oil collected from the industrial sector, e.g., cutting, stamping and coolant oils, which is collected on a direct basis or is collected from oily-water dehydrating facilities. This collected oil contains organo-metallic additives such as zinc dialkylthiophosphate from the original lubricating oil formulation, sludge formed in the engine, and water. The used oil may also contain contaminants such as waste grease, brake fluid, transmission oil, transformer oil, railroad lubricant, crude oil, antifreeze, dry cleaning fluid, degreasing solvents such as trichloroethylene, edible fats and oils, mineral acids, soot, earth and waste of unknown origin. Reclaiming of waste oil is largely carried out by small processors using various processes tailored to the available waste oil, product demands, and local environmental considerations. Such processes at a minimum include partial de-watering and coarse filtering. Some more sophisticated processors may practice chemical demetallizing or distillation. The presence of organo-metallics in waste oils such as zinc dialkylthiophosphate results in decomposition of the zinc dialkyldithiophospnate to form a carbonaceous layer rich in zinc and often other metals such as calcium, magnesium and other metals present as additives and thus difficult if not impossible to process. The carbonaceous layer containing the various metals forms rapidly on heated surfaces and can develop to a thickness of more than 1 mm in 24 hours. This layer not only reduces the heat transfer coefficient of tubular heaters rapidly, it also results in substantial or total occlusion of these tubes within a few days. Successful reclaiming processes require the reduction of the organo-metallics (or ash) content to a level at which the hot oil does not foul heated surfaces. Such reduction can be carried out by chemical processes which include reacting cation phosphate or cation sulfate with the chemically bonded metal to form metallic phosphate or metallic sulfate. U.S. Pat. No. 4,432,865 to Norman, the contents of which are incorporated herein by reference, discloses contacting used motor oil with polyfunctional mineral acid and polyhydroxy compound to react with undesired contaminants to form easily removable reaction products. These chemical processes suffer from attendant disposal problems depending on the metal by-products formed. Ash content can also be reduced by heating the used lubricating oil to decompose the organo-metallic additives. However, indirect heat exchange surfaces cannot be maintained above 400.degree. F. (204.degree. C.) for extended periods without extensive fouling and deposition of metals from the additives. Used lubricating oils can be heated to an additive decomposition temperature of 400.degree. F. (204.degree. C.) to 1000.degree. F. (538.degree. C.) by direct heat exchange by mixing with a heated product oil as disclosed in U.S. Pat. No. 5,447,628 to Harrison, et al., the contents of which are incorporated herein by reference. However, dilution of the product oil with used oil requires reprocessing already processed product oil . . . UOP's Hy-Lube process, described in U.S. Pat. Nos. 5,244,565 and 5,302,282, and many more, uses a hot circulating hydrogen stream as a heating medium to avoid deposition of decomposed organo-metallic compounds on heating surfaces. The problem of fouling of heated surfaces can be ameliorated to some extent by gentler heating. Some processes, such as the fixed bed version of catalytic cracking, the Houdry process, used a molten salt bath to provide controlled, somewhat gentle heating of vaporized liquid hydrocarbon passing through tubes of catalyst immersed in the salt bath. Molten metal baths have also been used as a convenient way to heat difficult to processes substances to a control temperature, e.g., flammability of some plastics is tested by putting a flask with plastic into a bath of molten metal. Use of molten salt bath, or molten metal bath, or condensing high temperature vapor, could be used to reduce uneven heating of heat exchange surface and thereby reduce dT across a metal surface and perhaps slow the fouling of metal surfaces in ULO service, but the additives in the ULO would still tend to decompose on the hottest surface, which would be the heat exchanger tubes. Although not related to ULO heating, in addition to the use of molten metal or molten salt for indirect heating as discussed above, there has been use, either commercial, or reported in the patent literature, of use of molten metal for direct contact heating of various substances. The float process for making glass is almost 50 years old. Molten metal, primarily lead, for heating coal or shale has been practiced in one form or another for almost 100 years. There are recent reports and patents on use of molten metal baths for waste pyrolysis, and conversion of latex, by heating ground up plants in a metal bath to make an oily overhead product. Also somewhat related, but even more different than anything discussed above, is the HyMelt® process, using molten iron beds for dissolution of various feed stocks. Temperatures in the HyMelt process are so high that if a liquid hydrocarbon feed is fed to a HyMelt reactor, the feed almost instantaneously dissociates in hydrogen and carbon, with the carbon dissolving in the molten iron. This is an excellent process for dissociating a hydrocarbon into its elemental constituents, but may be overkill for, e.g., reprocessing ULO, when all that is needed is enough heating to vaporize the lube boiling range components. Extensive work has also been done on use of molten salt baths to oxidize unwanted and difficult to process streams. Usually the salt baths are heat sources and reagents, i.e., intended to react with the feed, as reported in U.S. Pat. No. 3,845,910 or 4,602,574, which are incorporated by reference. Some researchers took the position that fouling of metal surfaces during ULO processing was going to happen, and that the best way to deal with it was to inject something into the ULO which would scrub the metal clean, i.e., injecting an abrasive material. Solvent extraction with light paraffin solvents such as propane, butane, pentane and mixtures thereof have been practiced by Interline and others. Details of the Interline Process are provided in U.S. Pat. No. 5,286,380 and U.S. Pat. No. 5,556,548. While the extraction approach seems like an elegant solution to the problem of processing ULO, the process may be relatively expensive to operate. Their quarterly report of May 15, 2002, reports that “It has become evident that demanding royalties based on production is impractical in many situations and countries. Unless and until the re-refined oil produced in a plant can be sold at prices comparable to base lubricating oils, collecting royalties based on production will be difficult. This reality was experienced in Korea, where the royalty was terminated for the first plant, and in England where the royalties were reduced and deferred until the plant becomes profitable. A breakthrough in ULO processing occurred with direct contact heating of the ULO with steam or a non-hydrogenating gas. This approach solved the problem of zinc additive decomposition fouling of hot metal surfaces, by ensuring that the metal surfaces holding the ULO were always relatively cool. The hottest spot in these ULO process was the point of vapor injection. Decomposing additives had only themselves to condense upon. Such a vapor injection ULO process was disclosed in my earlier patent, U.S. Pat. No. 6,068,759, Process for Recovering Lube Oil Base Stocks from Used Motor Oil . . . and in U.S. Pat. No. 6,447,672, Continuous Plural State Heated Vapor Injection Process for Recovering Lube Oil Base Stocks from Used Motor Oil . . . Other variations on the theme of ULO vapor injection processes are disclosed in U.S. Pat. No. 6,402,937 Pumped Recycle Vapor and U.S. Pat. No. 6,402,938, Vaporization of Used Motor Oil with Non-hydrogenating Recycle Vapor, which are incorporated by reference. The “state of the art” of used motor oil processing could be summarized as follows: Chemical additive and extraction approaches can be used to react with, or extract everything but, zinc additives, but costs associated with such processes are apparently high, as evidenced by little commercial use. Additives could be extracted, but the operating costs are high. Indirect heating, in a fired heater, causes rapid fouling of metal surfaces. Using milder heating, via a double boiler approach or molten metal heating medium, can minimize but not eliminate fouling on hot metal surfaces. Direct contact heating with high pressure hydrogen may eliminate fouling but requires high capital and operating expenses. Direct contact heating, with recycled product oil, helps but requires processing the ULO twice. Oxidation, either by burning as a low grade fuel, or perhaps as part of a salt bath oxidation process for waste streams. Direct contact heating with steam or non-hydrogenating vapor, as reported in my U.S. Pat. No. 6,068,759 and the related patents discussed above, is believed to be the best available technology. This approach requires only moderate capital investment and moderate operating expense when steam is the injected vapor, but the process can create a water disposal problem and is thermally less efficient because the latent heat of water is lost when the steam is condensed against cooling water or air in a heat exchanger. When other vapors are injected for heating e.g., propane, the water problem goes away but large volumes of vapor are needed to provide sufficient heat input, so costs increase to heat and recycle such vapor streams. Although my earlier work, steam injection for direct contact heat exchange, solved the worst problem, fouling on hot metal surfaces, it had some deficiencies as briefly noted above. I wanted an even better approach. I thought about steam injection. The steam injection process seemed nice and simple, because it was easy to heat water to make steam. Unfortunately, using large amounts of water created a potential water disposal problem and produced a relatively “wet” plant, with many potential areas for corrosion as the steam condensed. Re-using the condensed water was possible, but there are concerns about the amount of water treatment required to remove chlorides, etc, so that corrosion and/or plugging of the tubes in the fired heater would not be a problem. Large volumes of steam were required, which resulted in relatively large plant volumes, at least until some or all of the injected steam was condensed. I realized that although the use of steam was a great advance in the art, it might not be the best approach. The “pumped vapor” approach, use of propane or other recycle hydrocarbon vapor eliminates many concerns about water, but required a more complicated plant to recycle the hydrocarbon vapor. Large molar volumes of injected vapor are needed because of the relatively low heat capacity of hydrocarbon vapors. Condensation and separation of multiple hydrocarbon species, both injected heating vapors and recovered lubricating components, is more complicated than cooling everything and allowing water and oil to separate as separate phases. I wanted to retain the beneficial features of heating the ULO by injecting something hot into it, but avoid the problems created by using either steam or a light hydrocarbon vapor as the heating medium. I found a way to overcome these deficiencies, by using a non-pyrolizing molten salt bath as the heating fluid. There are many salts available which are fluid at relatively low temperatures which have ideal properties for use herein. They are relatively non-corrosive, especially when used in a reducing atmosphere. They are inexpensive and easy to contain. Molten salt is sufficiently dense to hold a lot of heat, permitting reasonably efficient heating of waste streams. They are not volatile, so they do not contribute to air or water pollution. They are immiscible with ULO so the decomposition products and trash found in the ULO can be easily removed from the molten salt bath. Molten salt also permits a flexible design approach, permitting injection of the molten salt into the oil or vice versa, though not necessarily with equivalent results. When oil is injected into a molten salt bath, it is easy to increase or decrease process severity by changing the depth of molten salt in the bath or the temperature of the salt or the pressure in the molten salt bath.
{ "pile_set_name": "USPTO Backgrounds" }
1. Field of the Invention This invention relates to transistorized electrical circuitry and more particularly to a protective circuit for protecting high power transistors operated in parallel and at high frequencies. 2. Description of the Prior Art In order to obtain high power at high frequencies, it is commonly necessary to parallel a plurality of transistors which are then operated as a composite group of like devices. This configuration can suffer catastrophic failure, however, when one of the transistors fails in a shorted condition because in a parallel circuit configuration all the bases are connected in parallel. The method heretofore used to protect the good transistors was to connect an electrical fuse in series with each base of all the transistors. Collector voltage (B+) was thus applied to all the bases when one of the transistors shorted until the fuse in the base of the shorted transistor blew. The elapsed time for clearing a fault, due to the nature of the fuse itself, was in the order of several milliseconds. However, due to the short length of "safe operating area" time for the remaining good transistors which is in the order of a few nanoseconds, the fuse will not operate sufficiently fast to prevent damage to the remaining transistors. The prior art protective schemes thus were adapted to protect the transistor to which the fuse was connected as opposed to protecting the balance of the remaining transistors from the failure of one transistor. The applicant has conducted a preliminary patentability search and has developed the following references which are hereby made of record: U.s. pat. No. 3,083,303, W.S. Knowles et al. PA1 U.s. pat. No. 3,490,031, I.R. Marcus et al. PA1 U.s. pat. No. 3,703,648, J.A. Wrabel. PA1 U.s. pat. No. 3,729,655, W. Gratzke.
{ "pile_set_name": "USPTO Backgrounds" }
1. Technical Field The present invention relates to an apparatus for testing bending strength, and particularly to an apparatus for testing hand-held device's bending strength. 2. General Background Hand-held devices, such as mobile phones, media players, etc, are easily damaged by inadvertent pressing stress. For example, if a user of the hand-held device sits down or squats down with the hand-held device is in the trousers' pocket of the user, the hand-held device will be subjected to the pressing stress. If the hand-held device is not strong enough, a liquid crystal display or an antenna of the hand-held device may be damaged. For reasons mentioned above, after a hand-held device prototype is created, it is necessary to perform a bending strength test on the hand-held device prototype before the hand-held device is mass produced. Therefore, what is needed is an apparatus for testing hand-held device' bending strength.
{ "pile_set_name": "USPTO Backgrounds" }
This invention relates to electrical signal observing devices in which an electrical signal to be measured is sampled with an optical signal to obtain an observation signal, and more particularly to an electrical signal observing device which has a simple construction and a high efficiency. Generally, repetitive high-speed electrical signals are measured with a sampling type oscilloscope (whose resolution is 20 to 30 pico-seconds at best) and single phenomena are measured with a real time oscilloscope (whose resolution is about 300 pico-seconds at best). However, the signals to be measured have increased in speed, and, consequently there is a demand for measurement of electrical signals with higher resolution. One way to obtain higher resolution is illustrated by the voltage measuring device of U.S. Pat. No. 4,446,425. In this device, an optical modulator comprising a non-linear optical medium (Pockels cell), polarizer, analyzer and compensator is used to sample an electrical signal with a short pulse light beam. However, using the technique of sampling an electrical signal with a short pulse light beam is not applicable to the measurement of single event phenomena. On the other hand, European Patent Application (OPI) No. 197196 (the term "OPI" as used herein means an "unexamined published application") discloses an electrically-electron optical oscilloscope in which, instead of a short pulse light beam, a continuous wave light beam (CW light beam) is employed. The above-described conventional device employs an optical modulator including a non-linear optical medium such as a Pockels cell which is generally expensive, and difficult to handle. Furthermore, this conventional device needs a polarizer, an analyzer and a compensator in addition to the non-linear optical medium, and therefore its optical system is unavoidably complicated and its adjustment is intricate and troublesome. The optical modulator using the non-linear optical medium merely changes light transmittance using the change in light polarization, and does not have amplification capabilities. Therefore this device cannot effectively utilize the incident light and has a low S/N ratio.
{ "pile_set_name": "USPTO Backgrounds" }
1. Field of the Invention This invention relates to swimming pools and, more particularly, to apparatus for sensing the height of water in a swimming pool and for controlling the flow of water to the swimming pool from a location remote from the swimming pool. 2. Description of the Prior Art There are several ways to control the height of water in a swimming pool. The simplest way to control the height of water in a swimming pool is by visual observation of the height of water in the pool, and upon visually determining that the water level is low, turning on a valve that allows water to flow to the swimming pool until the desired water level is observed. When the desired water level is observed, the water valve is manually closed, stopping the flow of water to the swimming pool. There have been other systems suggested for automatically controlling the height of water in a swimming pool, one such being suggested by U.S. Pat. No. 2,739,939. In the '939 patent, a float control operates a valve when the water level is low. When the water level is brought to a predetermined height, the float rises and turns off the valve. The system is not unlike the conventional float control valve used to control the height of water in a toilet tank. Another automatic system for controlling the height of water in a swimming pool is shown in U.S. Pat. No. 3,908,206. The swimming pool in the '206 patent is an above-ground swimming pool, and the water height control system is disposed adjacent to the swimming pool. A float is again used to control a valve which turns on and off a flow of water. In addition to the two patents discussed above, there have been other patents which include systems for automatically sensing the level of a liquid. Such liquid level patents include U.S. Pat. No. 352,647, U.S. Pat. No. 2,070,617, U.S. Pat. No. 2,305,102, and U.S. Pat. No. 4,014,010. While none of the four patents identified in the preceding paragraph are concerned with swimming pools, they each include means for sensing the height of a liquid in a tube of some type. In the '647 patent, a light source and a light sensitive element are disposed on opposite sides of a tube, with an opaque float disposed in the tube. The float is disposed between the light and the light sensitive element. When the float sinks with the level of the liquid in the tube, the resistance in an electrical circuit is changed and an alarm bell sounds. In the '617 patent, light sources and light sensitive elements are disposed on opposite sides of a tube in which the water level is indicated. An opaque float is also used in the tube. The changing water level causes the float to move, thus allowing light rays from the light sources to impinge upon light sensitive elements. Appropriate relays are used to control the flow of the water between the predetermined minimum and maximum heights. A light beam is used in conjunction with the '102 patent, but no float is used. In the '010 patent, a light source and a photo cell are disposed on opposite sides of a tube and an opaque float is disposed within the tube. Movement of the opaque float with respect to the light source and the photo cell controls an electrical circuit. When the float drops from between the light source and photo cell, an output to an alarm system which is activated, indicating that the fluid level has dropped below a predetermined minimum. In the two groups of patents, the patents concerned with swimming pools utilize floats to control the flow of water to a swimming pool. In the second group of patents, which have nothing to do with swimming pools, tubes are connected to a liquid tank of some kind, and a sensing system is used to determine the height or quantity of liquid (water) in the tank. In the '617 patent, the liquid sensing system also results in the actuation of a valve to keep the level in a boiler within predetermined limits. However, the other patents in the group are simply for reference purposes with respect to the height of the liquid and are not connected to a fill system. Moreover, all of the systems, in both groups of patents, are disposed adjacent to the liquid supply whose height they are concerned with. For swimming pools, it is advantageous to have the height sensing apparatus, particularly if it is associated with electricity, to be located remotely from the swimming pool so as to isolate the electrical system from the water system as much as possible. The apparatus of the present invention includes the remote sensing of the water height of a swimming pool and the electrical system(s) associated with the sensing and control system is thus isolated from the swimming pool.
{ "pile_set_name": "USPTO Backgrounds" }
To ensure device reliability, normally burn-in tests are conducted in burn-in testing systems once the devices are manufactured. A burn-in test is to apply an elevated voltage higher than the operating voltage to the control gate of a transistor of memory for a long period of time at an elevated temperature, usually 85° C. or above, making every element of the device subject to an overload, thereby exposing defects in the device in an early stage and detecting those devices with defects. A common burn-in testing system includes a burn-in testing apparatus and a Burn-In Board (BIB). For an improved yield, normally multiple devices under test are mounted on a large printed circuit board, i.e., the burn-in board. The devices under test on the burn-in board are connected in parallel, and their burn-in tests are performed at the same time. For more information on the burn-in testing apparatus and the structure of the burn-in board, please refer to Chinese Patent Application No. 200610163541.1 titled “Burn-in testing apparatus and burn-in testing board”. In the testing process, first a device under test is connected with a burn-in board, and the burn-in board is put into an environmental testing oven of a burn-in testing apparatus and is connected with a drive unit therein, then a testing environment desired for the test is achieved by adjusting conditions such as temperature according to a function to be tested for the device, and functional test is performed by the drive unit on the device through the burn-in board to detect the defect. At present, before carrying out the test, a burn-in testing apparatus manufacturer normally makes a burn-in board corresponding to a pin description of the device under test. When a new semiconductor device is to be tested, or when the pin description of the device is changed, a new burn-in board corresponding to the new pin description has to be made. Therefore, conventional burn-in testing systems have a low utilization of burn-in boards. Normally it takes weeks and thousands of dollars to make a burn-in board, and the low utilization may greatly increase production costs, prolong production cycle time, and hindering efficiency improving. In view of the foregoing problem, a burn-in board applicable to various types of semiconductor devices is desired in burn-in technologies for semiconductor devices.
{ "pile_set_name": "USPTO Backgrounds" }
This invention relates to a motor/generator provided with two or more rotors driven by a compound current. JP1999-275826A published by the Japanese Patent Office in 1999 discloses a motor/generator driven by a compound current. This motor/generator provides a cylindrical inner rotor and a hollow outer rotor disposed coaxially with a hollow cylindrical stator, there being a predetermined gap between the inside and outside of the stator. However, in this conventional motor/generator, even if linear voltage control is performed by an inverter, the maximum phase voltage of each phase can only be half of the alternating current voltage supplied from the inverter, and the utilization factor of the voltage is low. Moreover, as the phase voltage cannot be increased even if the current utility factor is improved by using a compound current, the required apparent power becomes small. If the alternating current voltage is increased by using a chopper, the alternating current voltage is high and the loss of the inverter increases. Furthermore, the number of power devices forming the inverter increased, the number of electronic parts, such as drivers, increased as a result, and this leads to restriction of layout flexibility and increased cost. It is therefore an object of this invention to provide a motor/generator which solves the above-mentioned problem. The details as well as other features and advantages of this invention are set forth in the remainder of the specification and are shown in the accompanying drawings. To achieve the above objects, this invention provides a motor/generator comprising a stator in which plural coils are disposed, rotors that are coaxially arranged with the stator and rotate by magnetic force generated by the coils, a device that generates a compound current from a direct current base on a command voltage to supply to the coils, and a controller functioning to calculate a neutral point potential and calculate the command voltage with reference to the neutral point potential.
{ "pile_set_name": "USPTO Backgrounds" }
This invention relates generally to perceiving rotation of an object. More specifically, the invention relates to the devices used for making such perceptions. In greater specificity the invention relates to an optical technique for detecting rotation or angular displacement of an object by utilizing the technology known as micro-electro-mechanical systems or xe2x80x9cMEMSxe2x80x9d. Typical MEMS-based gyroscopes use capacitive pick-offs to detect change in angular acceleration/rotation. This technique, however, is limited in sensitivity. This is particularly true as dimensions of the device are decreased, as capacitance varies linearly with area. One technique to increase the capacitance and thus sensitivity of these designs is to decrease the distance between the two parallel plates that form the capacitor. Decreasing this dimension presents its own problems, as variations in thickness and spacing across the surface of the parallel plates then play a much greater role in other performance shortcomings. Larger parallel plate spacings, on the order of a few microns, are therefore generally utilized to help increase repeatability and percent uniformity. Capacitive techniques result in a passive means of detection, meaning that additional low noise amplification and filtering circuitry must be employed to extract accurate rotational rates. Piezo-electric techniques cannot provide nearly enough sensitivity for many desired applications. Because of low steady-state signal levels, high sensitivity MEMS-based gyroscopes are difficult to realize. To gain the full potential of MEMS based gyroscopes, significant improvement in sensitivity over the prior-based methods must be made. Traditional gyroscope design incorporates a rapidly spinning element. The invention does not use a rapidly spinning component but instead uses a structure that allows the sensing of rotational movement and position of one structure with respect to another. The invention is based upon the integration of an optical resonant cavity and a photodiode. This combination is used to detect minute perturbations due to angular acceleration, such as those that may be generated by the Coriolis Force, perturbations due to constant angular acceleration otherwise known as rotational velocity, and angular displacement of one object with respect to another. For example, a Fabry-Perot cavity, consisting of two parallel semitransparent mirrors, can be used in conjunction with a light source. One of the two mirrors is fixed in position while the other is allowed to rotate with respect to the first mirror around a fixed rotational axis. A resonant cavity is thereby formed on either side of the axis, though only a single upper and lower mirror are used. If monochromatic light is used to illuminate the upper mirror of the cavity and the distance between the upper and lower mirrors is an integral multiple of half wavelengths of this light, then a resonant condition will exist and the transmission of light through the mirrors will be optimized. As the upper mirror is rotated, the distance between the two mirrors will become altered. One resonant cavity will xe2x80x9cseexe2x80x9d a decrease in cavity length while the other will see an equal but opposite increase in cavity length. As the distance between the two mirrors is changed, the light transmitting the mirrors will also be changed. Photodiodes integrated on either side of the torsional/rotational axis are used to sense the change in distance as a change in photo-generated current. By monitoring the change in photocurrent, the amount of change in rotation can be calculated. The photo-currents collected from the two cavities can be differentially amplified to further the sensitivity of the device. Other objects, advantages and new features of the invention will become apparent from the following detailed description of the invention when considered in conjunction with the accompanied drawings.
{ "pile_set_name": "USPTO Backgrounds" }
Advancements in communication and processing technologies have permitted the development, and implementation, of new types of communication systems. Generally, increased data rates of communication are permitted as a result of such advancements. Increasingly large amounts of data are permitted to be communicated within a selected time period. New types of communication services, requiring high data throughput rates are possible in such new types of communication systems. Multimedia communication services are exemplary of new types of communication services, permitted now due to such communication-technology advancements. And, many other types of communication services are similarly also possible in communication systems constructed to take advantage of the communication technology advancements. New-generation cellular, and other radio, communication systems are developed to provide for packet-based communications. Most multimedia communication services are predicated upon packet-based communications. Other communication services are analogously also migrating towards packet-based communications. Different types of communication schemes are used, or proposed for use, in such new communication systems. For instance, at least one WLAN (wireless local area network), to be operable pursuant to a variant of the IEEE (Institute of Electrical and Electronic Engineers) 802.11 specification proposes to utilize OFDM (orthogonal frequency division multiplexing) techniques. OFDM effectively forms a hybrid of a multi-carrier modulation (MCM) and frequency shift keying (FSK) modulation. Frequency-divided carrier frequencies are defined in an OFDM system such as the proposed, WLAN, and the carriers are selected to be orthogonal to one another, such as by separating the carriers by integer multiples of the inverses of symbol durations of parallel bit streams that are to be applied thereto. The orthogonal carriers are transmitted simultaneously, thereby permitting an entire allocated channel to be occupied through an aggregated sum of narrow, orthogonal subbands. The WLAN, operational pursuant to the variant of the IEEE 802.11 standard. Packet-based communication schemes generally utilize digital communication techniques, and communication stations operable in many packet-based communication systems require the utilization of modems (MOdulators/DEModulators). While operation of modems vary, depending upon the type of communication scheme in which the modems are to be operable, certain characteristics of the modems are important. Characteristics include the cost, size, complexity, programmability, scalability, simplicity, and modifiability of the modem. Existing modem constructions are sometimes categorized according to their general operational features. One modem-type category forms a data-flow style modem. Such a modem construction is formed of fixed circuits. The fixed circuits are interconnected by way of fixed connections, and the circuits are driven by a common, centrally-controlled timing circuit. The fixed circuits generally are operable in only one manner, e.g., a circuit is capable of executing only a single algorithm. While acceptable for many implementations, a conventional, data-flow construction is unable readily to be modified, is generally of relatively large physical dimensions, is relatively power-consumptive, and exhibits difficulty of design maintenance. Lack of programmability and scalability are disadvantageous from an economic viewpoint and further adaptability. Another modem-type category forms a DSP-style modem, i.e., a modem that utilizes a DSP (digital signal processor) in its construction. This type of modem is of improved programmability characteristics relative to a data-flow modem type. That is to say, the details of the modem are contained in software code and data arrangement, all storable in a memory device rather than in a particular interconnection of fixed circuits. However, the specialized nature of the DSP, and hardware structures ancillary thereto, limit the modifiability of the modem constructed therefrom. While more modifiable than a data-flow style modem construction, this modem type also lacks full programmability. A modem of this construction type also cannot therefore be modified easily. Execution of modem-function algorithms at a processor having a vector architecture would be advantageous as the speed at which the operations can be effectuated would be increased. Fast Fourier Transforming and Viterbi decoding in the demodulator part of the modem would particularly be facilitated. To date, however, a modem construction type has not been provided that fully takes advantage of the features of a processor having a vector architecture. If a modem utilizing such a processor could be provided modem operation would be facilitated and improved programmability and scalability to the modem would be provided. A modem construction type utilizing such a vector architecture would be particularly advantageously implemented in a communication system that utilizes OFDM. It is in light of this background information related to modem operation that the significant improvements of the present invention have evolved.
{ "pile_set_name": "USPTO Backgrounds" }
Millions of gallons of fuel oil and its equivalent are discarded every year through the disposal of plastic wastes and other waste material. Recycling of these wastes is of increasing importance as incineration and landfilling become more expensive and the acceptance of these methods is decreasing. It should be noted that rubber and plastic wastes are produced originally from crude oil and can be thermally cracked into fuels or petrochemicals. However, these wastes generally contain inorganic materials, fibers, glass, dust and poor thermal conducting materials, which are far more difficult to be treated effectively. The quantity of organic-containing solid wastes increases rapidly at the rate of millions of tons per year. These organic-containing solid wastes are equivalent to approximately thousands of billions kcal, which is a huge amount of thermal potential energy/heat and about half are from petroleum products. These solid wastes such as: printed circuit board wastes, rubber wastes, plastic wastes, scrap tire, organic wastes from auto shredder residues, oil sludge/sediment etc., are usually mixed with inorganic materials i.e., iron-wires, metal, fiber, wood, glass. Generally organic matter pollutes the solid wastes and this increases the difficulty of resource recovering treatment. In general, thermal treatment of wastes can recover energy and resources. These technologies include incineration, pyrolysis, oil liquefaction and gasification. Wastes incineration produces CO2 and H2O, but also produces some particulate, heavy metals, halides, SOx, and NOx. The accumulated pollutants have a negative impact on the environment. In addition the emission of PCCDS and PCDFS is also a serious problem. Under the condition of absence of oxygen, the macro organic compounds are cracked into smaller molecules and are recovered as the light hydrocarbons gases and light oil in the pyrolysis process. However, successful operation process for commercial purposes with direct pyrolysis are very few. This is due to engineering and operational problems such as (a) low heat transfer coefficient of the solid organic matter which affect the efficiency of the pyrolysis process; (b) high viscosity of products make the pyrolysis process more difficult; and (c) the pyrolysis products are not economically attractive, generally. Many direct pyrolysis processes have been reported to have technical or economic difficulties. Indeed, pyrolysis is complicated by the fact that the polymeric material wastes are poor conductors and degradation of these macromolecules requires considerable amount of energy. The liquefaction process involves treating solid wastes with hot waste lubricating oil at temperatures between about 435-800° F. (about 225-425° C. and below general pyrolysis temperatures). Basically, organic macromolecules are soluble in heavy oils only if they are cracked effectively. Above about 435° F. (about 225° C.), the C—C bonds of the polymeric matrix can be disrupted and dissolved into the oil. In a typical liquefaction process, solid wastes are liquefied in hot oil or recycled product heavy oils at relatively low temperatures. The liquefaction process comprises the main step of heat transfer by the hot oil to swell the structure of the highly polymeric organic material and lead to selective bond breaking. Therefore, the major products are oils, and can be separated easily from the mixtures of inorganic materials. The reaction temperature is usually less than about 750° F. (about 400° C.), which is much lower than that of any other known thermal treatment technologies. The need of the gas treatment equipment also becomes much less due to the lower quantity of the gas products. This proposed process can treat wastes mentioned above such as: printed circuit board wastes, scrap tire, plastic wastes, and other difficult-to-treat wastes as well as used motor oil. Hence the goal of treating several kinds of wastes simultaneously can be achieved. The liquefaction process usually takes place in a liquefaction reactor, where the solid wastes are dissolved and suspended into the hot oil or heavy recycled product oils. While the agitation requirements of this process have been demonstrated to be undemanding, the high viscosity of the mixture in the liquefaction reactor, the presence of non-soluble polymers and inorganic debris, and the tendency of any unmelted feed waste to float on the viscous fluid surface must all be accounted for. In small-scale liquefaction processes, such as on the pilot plant scale, liquefaction reactors are small enough that screens can be used to filter out the solid material from the liquefied waste. In large-scale commercial plants, screens are generally not used. Thus, there are concerns about the behavior of solid materials in those commercial plants. Additionally, unmelted waste small-scale units generally don't have the floating problems that commercial scale plants have. Thus commercial plant liquefaction reactors must have an agitation device that is better suited to drawing unmelted waste below the liquid surface, while minimizing the amount of gas drawn below the liquid surface and the depth that trapped gas is forced to.
{ "pile_set_name": "USPTO Backgrounds" }
1. Field of Invention This invention pertains to the field of pre-stressed pre-tensioned precast concrete slabs to be used for paving in areas subject to vehicular traffic. 2. Description of the Prior Art Prestressed concrete is a mode of construction that overcomes concrete's inherent weakness in tension. When concrete is prestressed using one of three means available, longer spans can be created as measured against ordinary reinforced concrete. Traditional reinforced concrete uses steel rebar or other reinforcement material disposed within the concrete to reinforce it. Typically a swimming pool bottom is made in this manner. Prestressed concrete employs cables or strands to provide a clamping load which produces a compressive stress that can balance the tension stress that the concrete member would otherwise exhibit due to a bending load. Pre-stressed concrete can be either pre-tensioned, or post tensioned. Pre-tensioned concrete is cast around already tensioned tendons. The concrete is poured around the pre-tensioned cables or tendons, and the concrete adheres to the tendons or cables as the concrete hardens during the curing process. When the tension is released from the tendons/cables this tension is transferred to the hardened concrete and compression by static friction thus creating concrete in compression. To achieve the pre-tensioning, anchor points are attached on opposite ends of the casting bed, between which, the tendons or wires are stretched in a straight line. When the tension is released, the tension is transferred to the hardened concrete unit by static friction. Post-tensioned concrete is the mode for applying compression after the pouring and curing in situ of the concrete. There are two modes of doing so, one is called Bonded and the other is called Un-bonded. In the bonded version, plastic, steel, or aluminum ducts, or tubes are laid out in a finite area, and the concrete is poured over and around the series of parallel tubes. Post tension cables are deployed through the tubes. Once the concrete hardens, the tendons are anchored at one end and tensioned at the other end using hydraulic jacks or rams that now react against the hardened concrete. After reviewing the design specification to confirm that adequate tension has been placed on the tendons, the jacks are removed such that the tension is now applied directly to the concrete member. The ducts or tubes are then grouted closed to protect the tendons from corrosion and decomposition. Concrete slabs prepared in this manner are usually used for bridges and house construction for slabs on grade in areas where the soil is expansive. In the unbonded system each individual cable has freedom of movement relative to the concrete at all times. Each individual tendon is coated with a grease, often lithium based, and perhaps molybdenum sulfide would work also. Then the tendons are covered by an extruded plastic sheathing. The tension transfer arises from the tendons being connected to anchors embedded in the perimeter of the cast concrete slab. While the generalized discussion of post-tensioning serves as an introduction to the topic, more information can be obtained from the Post-tensioning institute which in the year 2011 is located at 8601 North Black Canyon Highway in Phoenix, Ariz. Pre-stressed, Pre-tensioned concrete can not only be used for buildings, but is used today in Bridge work and the manufacture of roads. Pre-stressed paving slabs can be laid into position during off-peak hours on nights and weekends. This minimizes lane closures, which can cause huge traffic backups, especially on highly traveled interstate freeways. The big advantage of using pre-stressed concrete slabs, is the relative speed of placement on site, less cracking, and the ability to use relatively thinner and longer slabs. Longer slabs reduce the number of joints that must be maintained. Basically whereas standard construction can take weeks for a project, the same project can be carried out in days using pre-stressed, pre-cast concrete slabs. Numerous patents that relate to a method of forming, installing and a system for attaching prefabricated pavement slab to a subbase, and to the pavement slab itself have been issued to Peter J. Smith and said patents have been assigned to the Fort Miller Group, Inc. of Greenwich, N.Y. Some of these patents include: U.S. Pat. No. 6,709,792Issued Mar. 24, 2004U.S. Pat. No. 6,607,329Issued Aug. 19, 2003U.S. Pat. No. 6,899,489Issued May 31, 2005U.S. Pat. No. 6,962,462Issued Nov. 8, 2005U.S. Pat. No. 7,004,674Issued Feb. 28, 2006 andU.S. Pat. No. 7,467,776Issued Dec. 23, 2008 Another inventor in this technology is Alfred A. Yee, whose two patents are assigned to Kwik Slab, LLC of Honolulu, Hi. His patents are U.S. Pat. No. 7,134,805 which issued on Nov. 14, 2006 and the published application 2005/0220539. The Fort Miller Group product(s) are sold under the brand Super Slab, whereas the Yee products are sold under the brand Kwik Slab. It is believed that none of the aforementioned eight references singly or in combination disclose or render obvious the invention of this current patent application. The reason that this assertion can be made is that the invention of this patent application relates to an entirely new technique for pre-stressing, pre-tensioning concrete slabs in 2 directions, not just one direction as has been the case with the prior art techniques. As hinted above, the invention herein relates to a procedure for pre-stressing, pre-tensioning concrete slabs both longitudinally and transversely. The process further relates to the utilization of these bi-directionally pre-stressed, pre-tensioned slabs in the laying of roadways. In order to better understand this invention it is necessary to lay the foundation—no pun intended—of the general technique for making roadway sections. As noted above, pre-stressing can be accomplished by pre-tensioning or post tensioning. Pre-tensioning is done in the concrete casting bed, prior to the pouring of concrete, while post tensioning is done after concrete is poured and sufficiently hardened. Most concrete roadways are normally laid in up to 224 foot lengths between expansion joints. These sections are made of a plurality of slabs 12 feet wide and 8 foot long. These slabs can be connected by a variant of a tongue and groove connection or some other type of joint. The joints are then grouted or otherwise treated to form a complete section of concrete roadway. This means that in this 224 foot span there will be 28 grout joints. 8 feet long×28=224 feet. Generally pre-stressing in the concrete casting bed of a 12 foot length is carried out by pre-tensioning in the 12 foot direction prior to the pouring of the concrete and post-tensioning through the use of tendons or wires in a duct system after installation. But the pre-tensioning in the prior art techniques is in only one direction, longitudinally. The process of this invention is pre-stressing, pre-tensioning the concrete in both directions, longitudinally and laterally using a pre-tension technique longitudinally and laterally. Optional post-tensioning may also be applied. This allows for the preparation of longer slabs, potentially up to 60 feet in length, thereby minimizing the number of joints to be grouted and maintained in each roadway section, and thus speeding up the installation process. The invention accordingly comprises the apparatus (casting bed) and the device (dual direction, pre-stressed, pre-tensioned) concrete slab and the process of making the device, each of which possesses the features, properties, the selection of components which are amplified in the following detailed disclosure, and the scope of the application of which will be indicated in the appended claims.
{ "pile_set_name": "USPTO Backgrounds" }
Information is being produced in the world at a prodigious rate. For example, in the medical informatics field, the number of publications each year has more than doubled in the last decade. While such growth brings with it enormous possibilities, it is becoming increasingly difficult to stay current with advancements in a particular field. To decipher all of this information, numerous tools have been developed that help an interested party both understand the field and provide input into the field. For example, websites have been created that provide user generated, or society generated, encyclopedias. Such tools provide a starting point (and sometimes an ending point) for researching a particular area of interest. While these tools have proven very useful for synthesizing information, their features, and thus their utility, have been limited. For example, such tools typically do not inform readers where to go to find recently published articles on a topic and/or allow users to create multiple articles on the same topic.
{ "pile_set_name": "USPTO Backgrounds" }
Telephone systems allow users to conduct real time two-way voice communication. Traditional land-line based telephone systems connect one telephone set to another through one or more switching centers, operated by one or more telephone companies, over a land-line based telephone network. Traditionally, a telephone connection is based on a circuit switched network. Current telephone systems may also use a packet switched network for a telephone connection. A packet switched network is typical in a computer data environment. Recent developments in the field of Voice over Internet Protocol (VoIP) allow the delivery of voice information using the Internet Protocol (IP), in which voice information is packaged in a digital form in discrete packets rather than in the traditional circuit-committed protocols of the public switched telephone network (PSTN). Cellular networks allow a cellular phone to connect to a nearby cellular base station through an air interface for wireless access to a telephone network. Recent developments in wireless telephone systems allow not only voice communications but also data communications. For example, cellular phones can now receive and send short messages through a Short Message Service (SMS). Web pages can now be retrieved through wireless cellular links and displayed on cellular phones. Wireless Application Protocol (WAP) has been developed to overcome the constraints of relatively slow and intermittent nature of wireless links to access information similar or identical to World Wide Web. Telephone companies provide a number of convenient features, such as call forwarding. Call forwarding of a telephone system allows a user of a phone at a given phone number to dial a specific sequence on the phone to cause the telephone system to forward incoming calls addressed to the phone number to another specified phone number indicated by the dialed sequence. Telephone systems are frequently used in conducting business. Telephone numbers are typically provided in advertisements, web sites, directories, etc., as a type of contact information to reach businesses, experts, persons, etc. The Internet is becoming an advertisement media to reach globally populated web users. Advertisements can be included in a web page that is frequently visited by web users. Typically, the advertisements included in the web pages contain only a limited amount of information (e.g., a small paragraph, an icon, etc.). The advertisements contain links to the web sites that provide further detailed information. In certain arrangements, the advertisers pay the advertisements based on the number of visits directed to their web sites by the links of the advertisements. Performance based advertising generally refers to a type of advertising in which an advertiser pays only for a measurable event that is a direct result of an advertisement being viewed by a consumer. For example, paid inclusion advertising is a form of performance-based search advertising. With paid inclusion advertising, an advertisement is included within a search result page of a key word search. Each selection (“click”) of the advertisement from the results page is the measurable event for which the advertiser pays. In other words, payment by the advertiser is on a per click basis. Another form of performance-based advertising includes paid placement advertising. Paid placement advertising is similar to paid inclusion advertising in that payment is on a per click basis. However, with paid placement advertising an advertiser ranks a particular advertisement so that it appears or is placed at a particular spot, e.g., at the top of a search engine result page, thereby to increase the odds of the advertisement being selected. Both forms of performance-based advertising, i.e., paid placement and paid inclusion, suffer from the limitation that an advertiser or participant within a paid placement or paid inclusion advertising program is required to have a web presence, in the form of a web page. However, there are advertisers that either (a) do not have web pages, or (b) have web pages that are not effective at capturing the value of a web visitor, and are therefore unable, or unwilling, to participate in the traditional performance-based advertising, as described above.
{ "pile_set_name": "USPTO Backgrounds" }
Waste water treatment systems used in the industry generally include, but are not limited to, the following treatment processes: grit removal, fine screening, flow equalization and primary clarification. The typical treatment processes are dependent on the velocity at which the waste water is moving through the system. Waste water, however, is not produced continually by humans, but instead is created in batch type processes, such as showering, flushing a toilet or operating a washing machine. Such water consumptive activities are generally repetitive resulting in daily, weekly, monthly and yearly diurnal flow patterns for a specific waste water treatment system. Accordingly, the volume of waste water produced, and the velocity of that waste water through the treatment system varies significantly throughout the day. In the prior art, grit removal is generally performed in a grit chamber which is velocity sensitive. The most common methods to remove grit are by reducing the velocity of the influent flow so that the grit settles out, or utilizing a circular channel/tank. The circular channel/tank is a hydro-cyclone that causes the grit to settle in a sump, separating the organics from the grit so that they can move forward to the biological processes. The grit is then pumped out of the sump to a grit washer and then discharged to a dumpster for disposal at a landfill. Fine screening is typically accomplished by placing a screen in an influent channel. The influent channel must have a minimum velocity of 1.25 feet per second to keep solids from settling out in the channel and a maximum velocity of 3.0 feet per second to keep solids from being forced through the screen. Such a flow is difficult to achieve due to the large variation in diurnal and pumped flow patterns. Typical primary clarifiers are also velocity sensitive with the heavy solids going to the base of the clarifier where they are pumped to a digester, the floatable solids, grease and scum are trapped and skimmed off the surface and the neutral buoyant solids/clarified waste water exits the basin via an effluent weir. Primary clarifiers are typically large tanks designed for gravity settling and may include electrical drives, flights and chains, rack arms and paddles or suction tubes and sludge pumps. Flow equalization typically occurs in a separate tank. The flow at the waste water plant is subject to travel times in the collection system, collection system design and pump station sizing. In general, larger collection systems use pump stations to lift the waste water to the treatment facility. The pumps are typically placed on variable-frequency drives in an attempt to provide a consistent uniform flow. The system of variable-frequency drives and pumps, however, fails in low and high flow conditions. The pumps must be designed for peak hourly flows and have minimal turn down capabilities. Traditionally, waste water treatment plants have static bar racks or mechanically cleaned bar screens in channels at the entrance of the waste water into the treatment facility. These influent channels are typically constructed of concrete so as to last the life of the facility and are designed for specific waste water volumes, velocities (1 to 3 feet per second), and the insertion of specific screening and grit removal equipment. The social behavior of flushing solids that should go to landfill such as baby wipes, diapers, swizzle sticks, condoms, tampon applicators, etc. creates issues for the operation of the waste water treatment facility. Many of these solids are neutrally buoyant or will float in the waste water. Elongated solids align with the flow and pass or are forced through the bar racks or mechanical screens because of the high flow. The flat sheet solids such as diapers and baby wipes cover the bar racks or screens causing the liquid level in the channel to rise and enter a bypass channel. These solids often end up creating issues in the treatment plant such as fouling pumps, valves, diffusers, and membranes ultimately ending up in the digester or sludge holding tank. The increase in frequency and intensity of storm events producing exceptional precipitation combined with leaky sewage collection systems produces greater volumes of waste water delivered to the waste water treatment plant. Changes in societal behavior are not likely to occur. The cost to repair or replace the aged collection systems of developed nations is not fiscally achievable in the time frame needed. Therefore, the limited cross-sectional area of a channel requires an innovative approach to solve the above issues. The solution must be efficient in consideration of the goal to convert energy consumptive waste water treatment plants to sustainable resource recovery facilities where possible. To accomplish the above, the influent channels must be replaced with tanks. Waste water design engineers and manufacturers of screening equipment recognize that high velocities and screening are in conflict. Yet the use of channels at the head of the waste water treatment process is still taught to engineering students today.
{ "pile_set_name": "USPTO Backgrounds" }
1. Field of the Invention The present invention relates to an axis or a roller body on the basis of at least one first thermoplastic material which is processable in an injection molding process. This invention refers also to a method of producing this axis or this roller body, and the use of this axis or roller body, as well. 2. Description of the Related Art Conventional rollers include an axis produced of metal or alloys. This axis is bonded in a separate production step manually or mechanically to one or several rubber roller(s) or to one or several roller(s) of a thermoplastic material. This bond may be an adhesive bonding, a pressing on, an engaging, etc. These rubber rollers can, however, be directly vulcanised onto this axis. Before such rollers can be used in modern electronic appliances such as printers, labelling machines, photo copy apparatuses, fax machines, etc., an extensive surface treatment such as grinding, roughening, etc. of the roller(s) is necessary. Due to the various materials, the processing and post processing of such rollers is correspondingly expensive. In order to achieve the desired geometry at the metal axis, such as for instance the diameter of the shaft, grooves, driving elements, etc., extensive post-treatments are necessary. The natural large weight and also the possible electrostatic charging of the metal axis can have negative influences. In order to minimize the drawbacks and the restrictions in selecting the geometry at the metal axis, axes have been developed on the basis of plastic materials. Such conventional plastic material axes incorporate, however, still the drawback of the separate assembling of the roller of rubber or of a thermoplastic material, coupled with the above mentioned extensive surface treatment. In the JP-A-03-24926 a method of producing guide rollers for the use in video cassette apparatuses (video tape recorder, VTR) and in cassette tapes for VTR is described. Until the date of this Japanese invention such guide rollers were produced of only one single material. This had the following drawbacks: At the outer portions of such guide rollers recessed areas occur which necessitate a post-treatment by a machine. Furthermore, when selecting materials always a compromise between glideability and wear resistance had to be reached. According to JP-A-03-24926 these drawbacks are eliminated by the use of 2 different materials.
{ "pile_set_name": "USPTO Backgrounds" }
1. Field of the Invention Broadly, the present invention relates to a self-actuating and self-locking flow cutoff valve. It particularly relates to use of such a valve in a nuclear reactor of the type which utilizes a plurality of fluid supported absorber elements to provide for the safe shutdown of the reactor. 2. Prior Art There are numerous applications wherein there is a need for a self-actuating, self-locking flow cutoff valve. The need is particularly great in the case of nuclear reactors of the type which utilize a plurality of fluid supported neutron absorber elements to ensure the capability for a safe shutdown of the reactor. More particularly, heretofore nuclear reactors were typically shut down by control rods which were introduced through the top of the core and raised from or lowered into the core by mechanical means such as a motor which operates via clutch gears or the like. In an emergency, the clutch would be disengaged and the control rods allowed to fall into the core to shut down the reactor. Such a system had certain disadvantages. Specifically, there is a possibility that a mechanical device such as the clutch could not be disengaged or that some fault may have occurred which would distort the passage through which the control rods have to pass causing them to bind and preventing full insertion of the rods into the core. In such instance, it would not be possible to shut down the reactor. Accordingly, considerable interest has been generated in the use of a plurality of fluid supported neutron absorbing elements which would fall under the influence of gravity into the core in the event of a loss of fluid flow. Thus the reactor could be shut down by the simple expedient of shutting off the flow and further, in the event of an unforeseen loss of fluid flow, the reactor also would be shut down automatically. U.S. Pat. No. 3,228,847 suggests a reactor control system which includes a control assembly for controlling neutronic flux. The control assembly comprises an inner tube extending from a nonactive region of the reactor into the active region, and an outer tube surrounding the inner tube and spaced therefrom. The outer tube has a closed end and the inner tube has an open end adjacent and spaced from the closed end of the outer tube. Neutron absorbing particles are positioned between the inner and outer tube for movement along the tube under the force of flow. The neutron absorbing particles are moved out of the active region of the reactor by fluid flow and fall back into the active region under the influence of gravity when the flow is shut off. U.S. Pat. No. 3,257,286 suggests a ball-type control for a nuclear reactor. A number of elongated conduits are positioned in the nuclear reactor so that the first section of the conduit is located within the core and an adjoining second section is located exteriorly of the core. Each conduit holds a number of individual bodies, each of which contains a high neutron absorption cross-section material. The movement of the neutron absorber bodies within the conduits is achieved by providing a source of pressurized fluid available to each end of the conduit for selectively positioning the neutron absorber bodies within the first and second sections of the conduit. It is stated that a fission reactor can be started up, shut down, or reactivity controlled during reactor operations by varying the location of the absorber bodies. U.S. Pat. No. 3,347,747 discloses a control organization and method for a nuclear reactor. The reactor is provided with a number of laterally spaced vertical passageways in the region of the core and distributed throughout the area thereof. The passageways include a lower portion which extends generally throughout the height of the core and an upper portion which extends above the core into the reactor vessel. Positioned within and confined in each passageway is a movable means which contains a poison and which is movable from a lower position within the region of the core to an upper position in the passageway, where it is generally above the core. The poison-containing means is moved by gravity to its lower position and is moved from its lower to its upper position by means of a fluid which is directed upward in the passageway. In U.S. Pat. No. 4,076,583 there is disclosed another control method for a nuclear reactor which comprises a plurality of elongated conduits extending through and above the core of a reactor. A plurality of neutron absorber elements are located within the conduit, and during normal operation form a stacked bed in the portion of the conduit extending above the core. That section of the conduit in which the stacked bed is formed is provided with a fluid bypass means, it having been found that such bypass means ensures the capability of reliably maintaining all of the absorber elements in the stacked bed and out of the core during normal operations and further minimizing the pressure drop of fluid flowing through the stacked bed during normal operation. While all of the foregoing suggested techniques appear to offer advantages over reliance solely on a control rod system, there is still room for improvement. More particularly, in all of these systems where gravity is relied upon to cause the absorber elements to move into the core, any residual fluid flow, even though it may be below the minimum for safe operation of the reactor, acts to retard the fall of the absorber elements. For example, in the event of a complete power failure, the inertia of a centrifugal pump would be sufficient to continue providing some flow after the loss of power and after the flow rate of fluid had dropped below the point at which the reactor should be shut down. Thus, clearly it would be advantageous to have a self-actuating flow cutoff valve in the fluid stream such that once the fluid flow dropped below a predetermined point, the valve would automatically close and substantially reduce the time required for the neutron absorbing elements to fall into the reactor core and safely shut it down. Further, in the event that there might be some erratic flow or surge of pressure or flow subsequent to it having declined below the safe level, such valve advantageously would be self-locking to prevent an inadvertent startup of the reactor by a resumption of fluid flow.
{ "pile_set_name": "USPTO Backgrounds" }
An electric vehicle is a vehicle that includes an electric propulsion system. The electric propulsion system may include an electric motor and a battery. Hybrid vehicles may also include a combustion engine as well as a regenerative power system that transfers excess power from the combustion engine to the electric propulsion system. Electric vehicles may be charged by a charging station. The charging systems may be placed in parking garages, parking lots, or consumer homes. The electric vehicle may be electrically coupled to the charging station using a cord. Depending on the electrical input to the charging system, which may vary in amplitude and in number of phases, different charging stations may be capable of charging the electric vehicle in different amounts of time. Inductive charging systems may be developed in which no cord between the charging station and the electric vehicle is necessary. For example, inductive charging systems may be imbedded in the ground (e.g., concrete or asphalt) below a parking spot in a garage or parking lot. Electricity flowing through the inductive charging system may inductively charge a battery in the electric vehicle. Inductive charging system may also be imbedded in roadways and charge electric vehicles while the vehicles are traveling on the roadways. Charging systems in roadways are a scarce resource that can be managed in order to benefit a maximum number of electric vehicles.
{ "pile_set_name": "USPTO Backgrounds" }
Colour correction of digital colour images is required in a number of image processing contexts. One important environment is in digital imaging. It is known to produce colour sensors by introducing an alternating pattern of colour filters onto the array of individual sensor elements of a device. Alternatively, an image sensor may capture information relating to a plurality of different wavelengths of light at each point of the sensor. It is, however, difficult to construct colour filters for such sensors which exactly match the spectral characteristics of our eyes or which exactly match the primary colours used in computers to represent or display images. For this reason it is necessary for the captured images to be processed to transform the sensed colours to the desired colour system. These issues apply to other contexts in which colour correction is required, such as for images generated by imaging systems having three CCD sensors (one for each colour plane) for example, by flatbed colour scanners, or by other imaging systems in which the colour image is formed from sets of separate, registered images. Certain of these issues may apply to other contexts also such as printing, where colour correction is used to map from one colour space to that of the printer such as, for example, from a standard RGB (Red, Green, Blue additive primary) colour space to the RGB space of the printer (prior to the final transformation to the physical CMYK (Cyan, Magenta, Yellow, Black subtractive primary) colour space of the printer). When processing an image in order to transform it from one colour space to another, it is desirable to avoid mixing noise from a noisy channel, such as the blue colour channel for example, into a less noisy channel, such as the green colour channel for example. Co-pending United Kingdom Patent Application Number 0118456.3, incorporated herein by reference, discloses a method for the colour correction of images. An image to be processed is split into low and high frequency components and colour correction is applied to the low frequency component only. In this manner, the effect of noise is reduced during the colour correction process as the higher spatial frequency component of the image, which generally carries a larger proportion of the noise in an image, has no colour correction applied to it. The process of GB0118456.3 is suitable for modest transforms within the same basic colour space such as RGB to RGB, but it does not work particularly well in more extreme situations such as when transforming from complementary colours such as CMY to the primary RGB colours, for example. Both Japanese Patent Application No. 2003-110860 and “Suppression of Noise Amplification During Colour Correction”, Kharitonenko et al., IEEE Transactions on Consumer Electronics, Vol. 48, No. 2, May 2002 (Published), pp. 229-233 describe processes for colour correction of images. A further enhancement of GB0118456.3 is described in U.S. patent application Ser. No. 10/216,648. Therein, an adjustment may be applied to the high frequency image before recombining it with the colour corrected low frequency image in order to provide additional colour correction around areas of highly chromatic edges. Despite this improvement and the fact that any high frequency image component processing only occurs around highly chromatic edges, noise is still introduced into the final transformed image. Furthermore, the method of Ser. No. 10/216,648 is limited in its ability to convert an image from one colour space into a different one and is only suitable for transformation between broadly similar colour spaces.
{ "pile_set_name": "USPTO Backgrounds" }
The embodiments described herein relate generally to the field of ophthalmic therapies and more particularly to the use of a microneedle for delivery and/or removal of a substance, such as a fluid therapeutic agent into and/or from ocular tissues for treatment of the eye. Although needles are used in transdermal and intraocular drug delivery, there remains a need for improved microneedle devices and methods, particularly for delivery of substances (e.g., drugs) into the posterior region of the eye. Many inflammatory and proliferative diseases in the posterior region (or other regions) of the eye require long-term pharmacological treatment. Examples of such diseases include macular degeneration, diabetic retinopathy, and uveitis. It is often difficult to deliver effective doses of a drug to the back of the eye using conventional delivery methods such as topical application or an intravitreal administration (IVT), which has poor efficacy, and systemic administration, which often causes significant side effects. For example, while eye drops are useful in treating conditions affecting the exterior surface of the eye or tissues at the front of the eye, the eye drops are often not sufficiently conveyed to the back of the eye, as may be required for the treatment of some of the retinal diseases listed above. Although there have been advances in the past decade regarding the utilization of systemically delivered substances, there are obstacles to wide spread adoption of such methods. For example, in certain situations, direct injection into the eye (e.g., into the vitreous) using conventional 27 gauge or 30 gauge needles and syringes can be effective. Direct injection, however, can be associated with significant safety risks, and physicians often require professional training to effectively perform such methods. Moreover, in some instances, targeted injection of a therapeutic agent is desirable. In such instances, however, the relatively small anatomic structures of the eye often result in significant challenges to placing a needle at a target location using known devices and methods, especially as they pertain to placing the distal end of the needle at the desired depth within the eye. Furthermore, IVT administration can have side effects such as increased intraocular pressure or faster onset of cataract formation. In addition, many known methods of direct injection of a drug into the eye include inserting a needle or a cannula at an acute angle relative to a surface of the eye, which can make controlling the depth of insertion challenging. For example, some such methods include controlling the angular orientation of the needle such that the injected substance exits the needle at a particular location. Moreover, some known methods of injecting substances into ocular tissue include using complicated visualization system or sensors to control the placement of the needle or cannula. Known devices for ocular injection do not provide the mechanism for adjusting needle length so that the needle can be inserted into the eye to the desired depth. Known systems also do not provide a reliable mechanism for determining when the needle tip is in the desired location, for example, the suprachoroidal space (SCS) of the eye. Such shortcomings in known systems and methods are exacerbated because the size and thickness of various layers included in the eye can vary substantially from one person to another. For example, the thickness of the conjunctiva and the sclera can be substantially different and their true value cannot easily be predetermined via standard techniques. Furthermore, the thickness of these layers can also be different in different portions of the eye and at different times of the day in the same eye and location. Therefore, using known systems and methods it can be challenging to determine and/or adjust the length of the needle for puncturing the eye, such that a tip of the needle is at the desired depth, for example, the SCS. Too short a needle might not penetrate the sclera, and too long a needle can traverse beyond the SCS and damage the retina of the eye. Further, known systems do not provide a convenient way to detect the position of the needle tip within the eye. Because of the sensitivities associated with intraocular injection (e.g., the sensitivity of the tissue, the potential impact on intraocular pressure and the like), many known systems involve manual injection. More particularly, many known devices and methods include the user manually applying a force (e.g., via pushing a plunger with their thumb or fingers) to expel a fluid (e.g., a drug) into the eye. Because of the small needle size and/or the characteristics of the injected drug, some such devices and methods involve the use of force levels higher than that which users are comfortable with applying. For example, some studies have shown that users generally do not like to apply more than 2N force against the eye during ocular injection. Accordingly, in certain situations a user may not properly deliver the medicament using known systems and methods because of their reluctance to apply the force to fully expel the medicament. Moreover, injection into different target layers of the eye can cause variability in the amount of the force required for insertion of the needle and/or injection of the medicament. Different layers of the eye can have different densities. For example, the sclera generally has a higher density than the conjunctiva or the SCS. Differences in the density of the target region or layer can produce different backpressure against the needle exit, i.e., the tip of the needle from which the fluid emerges. Thus, injection into a relatively dense ocular material such as sclera requires more motive pressure to expel the medicament from the needle than is required when injecting a medicament into the SCS. Furthermore, the injection force to expel the medicament also depends on the density and viscosity of the liquid medicament, length of the needle, and diameter of the needle. To inject certain medicaments into the eye via desired needles (e.g., 27 gauge, 30 gauge, or even smaller) can require more force than many practitioners are comfortable applying. Intraocular injection can also lead to leakage of intraocular fluids (e.g., aqueous and vitreous humour) or the medicament from a delivery passageway formed by the needle penetrating into the ocular tissue. By way of example, if the medicament is delivered to the sclera instead of the target ocular tissue layer, for example, the SCS, the high backpressure of the sclera can force the medicament to leak from the insertion site. Known systems do not provide a convenient way to prevent leakage from insertion site, which can lead to discomfort and loss of medicament. This can prolong treatment as well as increase costs associated with the treatment. Thus, a need exists for improved devices and methods, which can assist in determining if the needle is at the correct depth, can facilitate injection of the medicament into ocular tissue, and/or can prevent leakage of ocular fluids and/or medicament form the insertion site.
{ "pile_set_name": "USPTO Backgrounds" }
calorimetry is a valuable tool in pharmaceutical research and development, providing information for decision making in drug lead discovery and optimization. Unlike the present high-throughput screening methods used by the industry, such as the affinity sensors, calorimetric analysis can provide very detailed information on the binding interaction between molecules. Calorimetry provides detailed thermodynamic information including the enthalpy and entropy of a reaction. The ability to measure enthalpy and determine entropy allows the drug development team to assess the relative contributions of enthalpy and entropy to a binding reaction. Enthalpy is driven by the number and type of bonds in the binding reaction. Entropy is driven by the geometry of the ligand and the binding site. Understanding the contributions of enthalpy and entropy is critical in drug development because it allows for the selection of compounds that are more readily optimized. Specifically, reactions that are enthalpy-driven tend to be favored due to their enhanced selectivity and reactivity. With current technology, initial high throughput screening and the first candidate drug selection is performed by affinity analysis. Only after the set of candidate drugs has been narrowed down to select few, candidates are analyzed by the two currently available calorimetry techniques, Differential Scanning calorimetry and Isothermal Titration calorimetry, to measure the thermochemical properties of a reaction. The limitations of the current generation of calorimeters include: Inadequate sensitivity for reactions with a low change in enthalpy Large amount of protein required (0.5 mg to 5 mg) Low experimental throughput because of both long experiment run times (60 to 90 minutes) and the need to sequentially run controls to assess the significance of the confounding effects. The potential confounding effects primarily include heating due to the mixing of dissimilar sample media (buffers with different pH, ionic strength, and solvents), the presence of DMSO from compound storage, and solvation. Furthermore, compounds with poor solubility frequently generate hits in high throughput screens. Unfortunately, the concentrations of these compounds required to meet the mass requirement for reagents (protein and its ligand) are often above the solubility limit. As a result, calorimetry studies on the interactions of these compounds with their targets cannot be done. Paradoxically, additional synthetic/medicinal chemistry is required before calorimetry can be used, but this chemistry work cannot be justified without the calorimetry data. The outcome of this is that potentially promising compounds are not pursued. The ability to analyze smaller amounts of reagents would reduce this need for concentration. Beyond pharmaceutical analysis, calorimetry is also valuable in many branches of materials science and chemistry. For example, calorimetry is useful for highly reactive or explosive compounds testing used in the design of chemical processes and safety equipment.
{ "pile_set_name": "USPTO Backgrounds" }