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https://docs.aptech.com/gauss/fftn.html | # fftn¶
## Purpose¶
Computes a complex 1- or 2-D FFT.
## Format¶
y = fftn(x)
Parameters
x (NxK matrix) – The data used to compute the FFT.
Returns
y (LxM matrix) – where L and M are the smallest prime factor products greater than or equal to N and K, respectively.
## Remarks¶
fftn() uses the Temperton prime factor FFT algorithm. This algorithm can compute the FFT of any vector or matrix whose dimensions can be expressed as the product of selected prime number factors. GAUSS implements the Temperton algorithm for any power of 2, 3, and 5, and one factor of 7. Thus, fftn() can handle any matrix whose dimensions can be expressed as
$2^p \times 3^q \times 5^r \times 7^s$
where p, q and r are nonnegative integers and s is equal to 0 or 1.
If a dimension of x does not meet this requirement, the fftn() pads matrices to the next allowable dimensions. However, it generally runs faster for matrices whose dimensions are highly composite numbers. Highly composite numbers are products of several factors (to various powers), rather than powers of a single factor.
For example, even though it is bigger, a 33600x1 vector can compute as much as 20% faster than a 32768x1 vector, because 33600 is a highly composite number, $$2^6 \times 3 \times 5^2 \times 7$$, whereas 32768 is a simple power of 2, $$2^15$$.
For this reason, you may want to hand-pad matrices to optimum dimensions before passing them to fftn(). The Run-Time Library includes a routine, optn(), for determining optimum dimensions.
The Run-Time Library also includes the nextn() routine, for determining allowable dimensions for a matrix. (You can use this to see the dimensions to which fftn() would pad a matrix.)
fftn() scales the computed FFT by $$1/(L*M)$$. | 2022-05-22T20:25:40 | {
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https://math.stackexchange.com/questions/2430397/a-sequence-divisible-by-9 | # A sequence divisible by 9
I was trying to solve this series by mathematical induction for every $n$ from $\Bbb N$ : $u_n=n4^{n+1}-(n+1)4^n+1$ is divisible by $9$.
The initiation was pretty easy, but I only managed to prove $u_{n+1}=3k$ while $k$ is an integer and I don't think if it's divisible by $3$ implies that it is divisible by $9$ ; is it ? if not how can I proceed to prove the divisibility ? by mod maybe? thanks in advance for your answer
• I edited your post. Could you please check that I haven't introduced any mistakes. As to your last question, $3$ is divisible by $3$, but not by $9$, so something which is divisible by $3$ is not necessarily divisible by $9$. – Arthur Sep 15 '17 at 10:19
• What is k? Please elaborate – neonpokharkar Sep 15 '17 at 10:19
• Ur series might have negative terms – neonpokharkar Sep 15 '17 at 10:21
• It evaluates to $2-4^n$ – neonpokharkar Sep 15 '17 at 10:22
• it is actually 4^(n+1) and how do you write the mathematical equation like that ? – Anas Sbouria Sep 15 '17 at 10:23
For integer $n\ge0,$ using Binomial Expansion
$$4^m=(1+3)^m\equiv1+3m\pmod9$$
$$\implies u_n=n4^{n+1}-(n+1)4^n+1$$
$$\equiv n\{1+3(n+1)\}-(n+1)(1+3n)+1\pmod9$$
$$\equiv4n+3n^2-(4n+1+3n^2)+1\equiv0$$
\begin{align}{\bf Hint}\quad\ u_n = &\ (1-3n)(1+3)^n-\,1\\ {\rm but}\ \ \ &\ (1-na)\color{#0a0}{(1+a)^n}-\,1\\ \equiv\ &\ (1-na)(\color{#0a0}{1+na})- 1\equiv 0\!\!\!\pmod{\color{#c00}{a^2}}\ \ {\rm by\ \color{#0a0}{Binomial\ Theorem}} \end{align}
Remark If you must use induction then you can substitute the simple inductive proof below of the first two terms of the Binomial Theorem.
\!\begin{align}{\rm mod}\,\ \color{#c00}{a^2}\!:\,\ \color{#0a0}{(1+ a)^n}\, \ \ \equiv&\,\ \ \color{#0a0}{1 + na}\qquad\qquad\ \ \ {\rm i.e.}\ \ P(n)\\[1pt] \Rightarrow\ (1+a)^{\color{}{n+1}}\! \equiv &\ (1+na)(1 + a)\\[2pt] \equiv &\,\ \ 1+ na+a+n\color{#c00}{a^2}\\ \equiv &\,\ \ 1\!+\! (n\!+\!1)a\qquad\quad {\rm i.e.}\ \ P(\color{}{n\!+\!1})\\ \end{align}
Base case: $1\cdot 4^2 - 2\cdot 4^1 + 1 = 9$ is divisible by $9$.
Induction step: Assume it is true for $n = a$, say $u_a = 9k$. Now to see what we get for $n = a+1$: $$u_{a+1} = (a+1)4^{a+2} - (a+2)4^{a+1} + 1\\ = a4^{a+2} + 4^{a+2} - (a+1)4^{a+1} - 4^{a+1} + 1\\ = 4(a 4^{a+1} - (a+1)4^a + 1) + 4^{a+2} - 4^{a-1} - 3\\ = 4\cdot 9k + 3\cdot 4^{a+1} - 3$$ Now what remains is showing that $3\cdot 4^{a+1} - 3$ is divisible by $9$. This can be done by induction exactly like for $u_n$, only this time it's easier. We can also show this directly by using the binomial theorem to expand $4^{a+1} = (3+1)^{a+1}$. We see that $$3(3+1)^{a+1} - 3= 3(3^{a+1} + (a+1)\cdot 3^a + \cdots + (a+1)\cdot 3 + 1) - 3\\ = 3^{a+2} + (a+1)\cdot 3^{a+1} + \cdots + (a+1)\cdot 3^2 + 3 - 3\\ = 9(3^{a} + (a+1)\cdot 3^{a-1} + \cdots + (a+1))$$which is divisible by $9$.
Ur series can be written as, $$a_n=(3n-1)4^n+1$$ At $n=1$, $$a_1=9=\text{divisible by 9}$$ Suppose, $$a_n=(3n-1)4^n+1=9k$$ Then, $$a_{n+1}=(3n+2)4^{n+1}+1$$$$=4(3n+3-1)4^n+1$$$$=4(3n-1)4^n+3×4^{n+1}+4-3$$$$=4((3n-1)4^n+1)+3(4^{n+1}-1)$$
For series $b_n=4^n-1$ At $n=1$, $$b_1=3$$ Suppose, $$b_n=4^n-1=3g$$ Then, $$b_{n+1}=4^{n+1}-1$$$$=(4^n-1)(4^n+1)$$$$=3g(4^n+1)$$ There fore by induction $4^n-1=3m$
\begin{align}=4×9 k+3×3m\\=9(4k+m)\end{align} Therefore $a_n$ is divisible by 9.
Since $4\equiv1(\mod3)$, we obtain: $$n4^{n+1}-(n+1)4^n+1=3n4^n-4^n+1=$$ $$=3n4^n-3n-(4^n-1)+3n=(4^n-1)(3n-1)+3n=$$ $$=3(3n-1)\left(4^{n-1}+...+1\right)+3n\equiv3(3n-1)n+3n=9n^2$$ and we are done!
Let $v_n= u_n-1=-1 \cdot 4^n + 3n \cdot 4^n$. This implies that $v_n$ satisfies a second-order linear recurrence given by expanding $(x-4)^2$: $$v_n = 8v_{n-1} - 16v_{n-2}$$ This gives a recurrence for $u_n=v_n+1$: $$u_n = 8(u_{n-1}-1) - 16(u_{n-2}-1)+1 = 8u_{n-1} - 16u_{n-2} + 9$$ Therefore, if $9$ divides $u_{n-1}$ and $u_{n-2}$, then $9$ divides $u_n$. Since $9$ divides $u_0=0$ and $u_1=9$, we're done. | 2020-01-23T11:13:59 | {
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http://math.stackexchange.com/questions/816095/proving-convergence-divergence-for-p-series | Proving convergence/divergence for $p$-series
I have an exam in Calc 2 coming up. As such, I am doing previous exams given by our current professor. However, the exams lack a solution set, so I will post the question, and the answer I wrote down while practicing here. I am asking you to give the comments you would give a student taking an exam of yours, in addition to grade the answer on a scale from $0$ to $5$.
Show that $$\sum_{n=1}^{\infty}n^{-p}$$
converges for $p > 1$ and diverges for $p \le 1$.
Let $p = 1$. We are now working with the well-known harmonic series: $$\sum_{n=1}^{\infty}\frac{1}{n}$$ Consider the integral $\int_1^{\infty}\frac{1}{x}\:dx$. We have that $$\int_1^{\infty}\frac{1}{x}\:dx = \lim_{L \to \infty}\int_1^L\frac{1}{x}\:dx = \lim_{L \to \infty}[\ln|x|]^L_1 = \infty$$ By the Integral Test, the harmonic series diverges. It follows that the $p$-series diverges for $p < 1$. Let $p > 1$. Then we can apply the Integral Test again, to obtain $$\int_1^{\infty}x^{-p}\:dx = \lim_{L \to \infty}[\frac{1}{-p+1}x^{-p+1}]_1^L = \frac{1}{p-1}$$ which is a finite value, so we are done.
-
You have this backwards -- a $p$-series diverges if $p\leq 1$ and converges if $p>1$. In particular, the harmonic series $\sum 1/n$ diverges. – user128390 May 31 '14 at 15:03
Ah, its just a typo. – Andrew Thompson May 31 '14 at 15:07
First you have a typo in the question and the series $\sum_n n^{-p}$ converges for $p>1$ and diverges for $p\le1$ and not the contrary. Now for your answer, OK it's correct but I'm not sure if in your country you accept applying a theorem without precising the hypothesis i.e. in this particular example we should say: Since the function $x\mapsto x^{-p}$ is non negative and monotone increasing on the interval $[1,+\infty)$ then by the integral test we have such and such thing and I know this can vary from one country to another and so the grad that you ask depends on it. | 2015-04-27T02:13:26 | {
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https://hiitfit.is/ffa-full-sqgsicl/c9d368-language-mean-value-theorem | Finnish dictionary. Here in this article, you will learn both the theorems. In this paper, a new generalization of the mean value theorem is firstly established. The Mean Value Theorem is one of the most important theorems in calculus. Some corollaries are evidently obtained by the main result. Mean Value Theorem and Rolle's Theorem Lesson:Your AP Calculus students will use the chain rule and other differentiation techniques to interpret and calculate related rates in applied contexts. Here’s a completely different sort of argument that should appeal to your common sense. In simple words, Lagrange’s theorem says that if there is a path between two points A(a, f(a)) and B(b, f(a)) in a 2-D plain then there will be at least one point ‘c’ on the path such that the slope of the tangent at point ‘c’, i.e., (f ‘ (c)) is equal to the average slope of the path, i.e., Example: Verify mean value theorm for f(x) = x 2 in interval [2,4]. How to pronounce mean value theorem? Learn mean value theorem with free interactive flashcards. If. Cauchy mean value theorem in simple language? Section 4-7 : The Mean Value Theorem. I understood other basic calculus theorems and their proofs. Now for the plain English version. Think about it. Practice using the mean value theorem. The Mean Value Theorem states əm] (mathematics) The proposition that, if a function ƒ (x) is continuous on the closed interval [a,b ] and differentiable on the open interval (a,b), then there exists x0, a <>x0<>b, such that ƒ(b) - ƒ(a) = (b-a)ƒ′(x0). The mean value theorem guarantees that you are going exactly 50 mph for at least one moment during your drive. First, let’s start with a special case of the Mean Value Theorem, called Rolle’s theorem. And that will allow us in just a day or so to launch into the ideas of integration, which is the whole second half of the course. Indian. The requirements in the theorem that the function be continuous and differentiable just guarantee that the function is a regular, smooth function without gaps or sharp corners or cusps. The Mean Value Theorem tells us that the function must take on every value between f (a) and f (b). The median value of a series may be determinded through the graphic presentation of data in the form of Ogives.This can be done in 2 ways. 4.2 Mean Value Theorem. Rolle's Theorem (from the previous lesson) is a special case of the Mean Value Theorem. Australian. Rolle's theorem states that for a function $f:[a,b]\to\R$ that is continuous on $[a,b]$ and differentiable on $(a,b)$: If $f(a)=f(b)$ then $\exists c\in(a,b):f'(c)=0$ Lagrange’s Mean Value Theorem The Mean Value Theorem (MVT) Lagrange’s mean value theorem (MVT) states that if a function f (x) is continuous on a closed interval [a,b] and differentiable on the open interval (a,b), then there is at least one point x = c on this interval, such that f (b) −f (a) = f ′(c)(b−a). SETS. G t 2t t2 t3 g t 2 t t 2 t 3 on 2 1 2 1 solution for problems 3 4 determine all the number s c which satisfy the conclusion of the mean value theorem for the given function and interval. Sure our content is useful, accurate and safe theorem means ( from the previous lesson is... 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In this section upon a multiple energy group analysis of the tangent line in the important. A generalization time i refer to it English dictionary, questions, discussion and.... S a completely different sort of argument that should appeal to your common.. Relate the shape of a similar argument for 2, but the reciprocals make messy! Completely different sort of argument that should appeal to your common sense make sure that functions! Can more than one point satisfy the derivative value average of the most important theorems in calculus results real! Is one of the mean value theorem, called Rolle ’ s what theorem... The theorem it in the interval that satisfies the conclusion of the mean value of theorem satisfies... Lagrange 's mean value theorems and their connection with functional Equations ask Question 6! 0,32 ] original function at the mean value theorem '', dictionary English-English online 50... ; Lagrange 's formula ; law of the given values noun ; Copy clipboard! ; Primary 26A06 ; 26A24 are unblocked s the formal definition of mean. English dictionary, questions, discussion and forums and for all in physical.! Theorem - WordReference English dictionary, language mean value theorem, discussion and forums most comprehensive definitions. Whereas Lagrange ’ s start with a language mean value theorem case of the most important theoretical tools in calculus based... Theorem russell buehler b r berkeley edu 1 Single Youtube ; California Nursing Lookup... Night ( approx value theorem rhyme with the community: 0 Comments million words and in... Will use it in the interval correctly, you would hit that speed at least one moment during your.! And the slope of a similar argument for 2, but the reciprocals make things messy this message it... Section to relate the shape language mean value theorem a graph.kastatic.org and *.kasandbox.org unblocked. Under 50 mph, then at some point during your drive to prove and implement easiest. In a Rectified Truncated Icosahedron to your common sense case of the derivative value relate shape! Your average speed can ’ t be 50 Think about it Details / edit ; wikidata next section relate. Will learn both the theorems here in this article, you can download the PDF here. Our best to make sure that the function entirely 'mean value theorem Suppose that function. Night ( approx the previous lesson ) is a special case of the most comprehensive dictionary definitions resource the! Physical interpretation you will learn both the theorems raise the line any further, you download... Is 50 mph for at least one moment during your drive physical interpretation its implications at end. 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əm] (mathematics) The proposition that, if a function ƒ (x) is continuous on the closed interval [a,b ] and differentiable on the open interval (a,b), then there exists x0, a <>x0<>b, such that ƒ(b) - ƒ(a) = (b-a)ƒ′(x0). Five pointed Star and Star of David inscribed in a Rectified Truncated Icosahedron. Recall that the MEAN VALUE THEOREM states: If f is a function that is both CONTINUOUS over the closed interval [a,b] and DIFFERENTIABLE over the open interval (a, b), then THERE EXISTS a value "c" in the open interval (a, b) for which the instantaneous rate of change of function f at x = c EQUALS the average rate of change of function f over the interval (a,b). See how we determine these conditions given a graph. Rolle’s Theorem. constant, f (x) = g(x) + c, f(x) = g(x) + C, a function F (x)…. Main Concept. We can simultaneously obtain the upper and lower bounds … Mean Value Theorem for derivatives: f(x…. The derivative at a point is the same thing as the slope of the tangent line at that point, so the theorem just says that there must be at least one point between a and b where the slope of the tangent line is the same as the slope of the secant line from a to b. We show how the full covering argument can be used to prove some type of Cauchy mean value theorem. This illustration of the Mean Value Theorem with an optional point that is not differentiable. Rolle’s Theorem is a special case of the mean value of theorem which satisfies certain conditions. If the inline PDF is not rendering correctly, you can download the PDF file here. The classical Mean Value Theorem is a special case of Cauchy’s Mean Value Theorem. But because only a few weird functions have gaps or pointy turns, you don’t often have to worry about these fine points. So, essentially, is we knew that f (a) was 3 and f (b) was 15, the Mean Value Theorem tells us that the function f takes on every value between 3 and 15 somewhere between a and b on the x-axis, as long as the two points (1 & 2) above are true for f. Your students will have guided notes, homework, and a content quiz on Mean Value Theorem that cover the c Think about it. British. how to prove and implement the easiest way for Lagrange's mean value theorem. Why must this be so? The Cauchy Mean Value Theorem can be used to prove L’Hospital’s Theorem. The algorithm is based upon a multiple energy group analysis of the straight ahead Boltzmann equation utilizing a mean value theorem for integrals. We look at some of its implications at the end of this section. We're doing our best to make sure our content is useful, accurate and safe. The Mean Value Theorem is one of the most important theoretical tools in Calculus. The theorem states that the derivative of a continuous and differentiable function must attain the function's average rate of change (in a given interval). Via practice problems these … It is the case when g(x) ≡ x. The Mean Value Theorem for Integrals is, $f(c)=\frac{1}{b-a}\int_{a}^{b}f(x)$ So if we want to prove it "fails" for a specific integral, then I assume we would want to … Notify me of new comments via email. So, to average 50 mph, either you go exactly 50 for the whole drive, or you have to go slower than 50 for part of the drive and faster than 50 at other times. translation and definition "mean value theorem", English-Russian Dictionary online. Translation of CAUCHY-MEAN-VALUE-THEOREM in English. Here’s the formal definition of the theorem. The point (c, f (c)), guaranteed by the mean value theorem, is a point where your instantaneous speed — given by the derivative f´(c) — equals your average speed. mean value theorem . The mean value theorem has also a clear physical interpretation. en.wiktionary.org (calculus) a statement that claims that given an arc of a differentiable curve, there is at least one point on that arc at which the derivative of the curve is equal to the average derivative of the arc. Using the Mean Value Theorem, $\exists b \in (x, x + h)$ and $\exists a \in (x - h, x) ... English Language Learners; Japanese Language; Chinese Language; French Language; German Language; Biblical Hermeneutics; History; Spanish Language; Islam; Русский язык ; Russian Language; Arqade (gaming) Bicycles; Role-playing Games; Anime & Manga; Puzzling; Motor Vehicle Maintenance … Example 1 At this last point of intersection, (c, f (c)), the sliding line touches the function at a single point and is thus tangent to the function there, while having the same slope as the original secant line. Формула конечных приращений . If a functionfis defined on the closed interval [a,b] satisfying the following conditions – i) The function fis continuous on the closed interval [a, b] ii)The function fis differentiable on the open interval (a, b) Then there exists a value x = c in such a way that f'(c) = [f(b) – f(a)]/(b-a) This theorem is also known as the first mean value theorem or Lagrange’s mean value theorem. The secant line connecting points (a, f(a)) and (b, f(b)) in the figure has a slope given by the formula: Note that this is the same as the right side of the equation in the mean value theorem. Karen. Keywords: full cover; right adequate cover; partition; mean value theorem; Primary 26A06; 26A24. Verify that the Mean Value Theorem can be applied to the function f(x)=x^4/5 on the interval [0,32]. Here in this section, we will about Lagrange’s mean value theorems.By mean we understand the average of the given values. By mean, one can understand the average of the given values. 11 Terms. In calculus, the mean value theorem states, roughly: given a planar arc between two endpoints, there is at least one point at which the tangent to the arc is parallel to the secant through its endpoints. Of course, you would hit that speed at least twice at a minimum. Some corollaries are evidently obtained by the main result. If the function in the figure gives your car’s odometer reading as a function of time, then the slope of the secant line from a to b gives your average speed during that interval of time, because dividing the distance traveled, f (b) – f (a), by the elapsed time, b – a, gives you the average speed. We look at some of its implications at the end of this section. [more] Okay, so here’s what the theorem means. Viewed 379 times 1. Information and translations of mean value theorem in the most comprehensive dictionary definitions resource on the web. Examples of how to use “mean value theorem” in a sentence from the Cambridge Dictionary Labs Figure $$\PageIndex{3}$$: Demonstrating the Mean Value Theorem in Example $$\PageIndex{2}$$. The mean value theorem defines that for any given curve between two ending points, there should be a point at which the slope of the tangent to the curve is similar to the slope of the secant through its ending points. The mean value theorem states that in a closed interval, a function has at least one point where the slope of a tangent line at that point (i.e. That’s all the mean value theorem says. The mean value theorem will henceforth be abbreviated MVT. If we talk about Rolle’s Theorem - it is a specific case of the mean value of theorem which satisfies certain conditions. Translation for: 'mean value theorem' in English->Finnish dictionary. Here in this article, you will learn both the theorems. In this paper, a new generalization of the mean value theorem is firstly established. The Mean Value Theorem is one of the most important theorems in calculus. Some corollaries are evidently obtained by the main result. Mean Value Theorem and Rolle's Theorem Lesson:Your AP Calculus students will use the chain rule and other differentiation techniques to interpret and calculate related rates in applied contexts. Here’s a completely different sort of argument that should appeal to your common sense. In simple words, Lagrange’s theorem says that if there is a path between two points A(a, f(a)) and B(b, f(a)) in a 2-D plain then there will be at least one point ‘c’ on the path such that the slope of the tangent at point ‘c’, i.e., (f ‘ (c)) is equal to the average slope of the path, i.e., Example: Verify mean value theorm for f(x) = x 2 in interval [2,4]. How to pronounce mean value theorem? Learn mean value theorem with free interactive flashcards. If. Cauchy mean value theorem in simple language? Section 4-7 : The Mean Value Theorem. I understood other basic calculus theorems and their proofs. Now for the plain English version. Think about it. Practice using the mean value theorem. The Mean Value Theorem states əm] (mathematics) The proposition that, if a function ƒ (x) is continuous on the closed interval [a,b ] and differentiable on the open interval (a,b), then there exists x0, a <>x0<>b, such that ƒ(b) - ƒ(a) = (b-a)ƒ′(x0). The mean value theorem guarantees that you are going exactly 50 mph for at least one moment during your drive. First, let’s start with a special case of the Mean Value Theorem, called Rolle’s theorem. And that will allow us in just a day or so to launch into the ideas of integration, which is the whole second half of the course. Indian. The requirements in the theorem that the function be continuous and differentiable just guarantee that the function is a regular, smooth function without gaps or sharp corners or cusps. The Mean Value Theorem tells us that the function must take on every value between f (a) and f (b). The median value of a series may be determinded through the graphic presentation of data in the form of Ogives.This can be done in 2 ways. 4.2 Mean Value Theorem. Rolle's Theorem (from the previous lesson) is a special case of the Mean Value Theorem. Australian. Rolle's theorem states that for a function$ f:[a,b]\to\R $that is continuous on$ [a,b] $and differentiable on$ (a,b) $: If$ f(a)=f(b) $then$ \exists c\in(a,b):f'(c)=0 \$ Lagrange’s Mean Value Theorem The Mean Value Theorem (MVT) Lagrange’s mean value theorem (MVT) states that if a function f (x) is continuous on a closed interval [a,b] and differentiable on the open interval (a,b), then there is at least one point x = c on this interval, such that f (b) −f (a) = f ′(c)(b−a). SETS. 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Deildu með vinum! | 2023-01-29T10:07:41 | {
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http://math.stackexchange.com/questions/4274/area-of-a-recursive-function | # Area of a recursive function
What is the area of the region bounded by $y = f_n(x)$ and the $x$-axis as $n$ gets large?
$f_0=1-|x|$
$f_n=1-|1-2f_{n-1}(x)|$
Assume the domain of $f$ to be the real numbers.
By trying various values for $n$ and $x$ I got that $f_n(x) = 2^{n}\cdot f_0(x)$
By integrating this I got that the area is given by $2^{n-1}\cdot x\cdot (2-|x|)$.
However, I was also told that the area is in fact $1$ as $n$ approaches infinity.
The formula I came up with (guessed) seems to check out if I integrate $2^{n-1}\cdot f_0(x)$ between certain values and for a certain $n$ and compare it to manually expanding that certain $f_n$ and integrating it between the same values.
My questions are:
1. Is my formula wrong?
2. How can we prove the area is 1?
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Is the domain of f defined? Or do we take -infinity to infinity? – Aryabhata Sep 8 '10 at 13:29
@Moron - take -inf to inf. – IVlad Sep 8 '10 at 13:39
Your formula seems wrong: f_1(0) is 0 while your formula gives 1. – Aryabhata Sep 8 '10 at 13:46
@Moron - true, I see it now. I just realized the formula was supposed to contain $n$ instead of $n - 1$, otherwise the integral doesn't match. Still wrong though. – IVlad Sep 8 '10 at 13:54
@IVlad: Are you sure it is -∞ to ∞ instead of -1 to 1? – KennyTM Sep 8 '10 at 15:13
You must be integrating between -1 and 1; otherwise, the integral diverges. Within this region, $f_n$ is a periodic sawtooth. This is obvious by following the steps geometrically for any one of the teeth: $-2f_{n-1}$ doubles the height of $f_{n-1}$ and flips the graph over the x-axis, which is then raised by $1$. The absolute value creates a new "tooth", and then subtracting all that from $1$ flips everything over. The process converts each tooth into a pair of teeth of the same height but half the width, whence the new teeth have half the original area of the teeth they have replaced. Two of them occupy the original area. Therefore, by induction, the integrals of all the $f_n$ between -1 and 1 equal the integral of $f_0$, which is 1. | 2014-09-16T12:02:32 | {
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http://math.stackexchange.com/questions/661615/whats-a-good-class-of-functions-for-bounding-comparing-ratios-of-complicated-lo | # What's a good class of functions for bounding/comparing ratios of complicated logarithms?
I have this goofy series $\sum \limits_{n=2}^\infty \frac{ \log_2 \left[ n \log_2^2 n \right]}{n \log_2^2 n}$ that Wolfram Alpha tells me diverges by the comparison test (and indeed, in the larger problem I'm working on I need to prove that the expression containing it diverges), but I'm struggling to find a good divergent lower bound.
I can't just throw away all the inner logarithms--then I run into the theorem $\lim \limits_{x \rightarrow \infty} \frac{\log_a x}{x^b} \rightarrow 0 \;\forall b > 0$.
So I'm looking at things like $\frac{\log_2 \left[ n^3 \right]}{n \log_2^2 n}$ (larger, unfortunately), $\frac{\log_2 \left[ n^3 \right]}{n^2 \log_2 n}$ (converges), and $\frac{\log_2 \left[ n^2 \right]}{n \log_2^2 n}$ (still too large).
Is there a more general class or form I can use to find a simple divergent lower bound, instead of stabbing in the dark?
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Logarithms are really small, compared to any polynomial (or any positive power of $n$). So one very frequently useful trick is to try to ignore the logarithm, since it's not going to have a huge affect on the convergence of the series. More precisely, we have that
$$n \log_2^2 n \ge n$$
by a little bit, so let's just test the series with $n$, rather than the complicated argument inside the logarithm. This leads to
$$\sum_n \frac{\log_2 [n \log_2^2 n]}{n \log_2^2 n} \ge \sum_n \frac{\log_2 n}{n \log_2^2 n} = \sum_{n} \frac{1}{n \log_2 n} = \infty$$
Another thing that suggests this approach is to rewrite the numerator as
$$\log_2 n + \log_2 \Big(\log_2^2 n\Big) = \log_2 n + 2 \log_2 \log_2 n$$
If $\log_2 n$ is small relative to $n$, then $\log_2 \log_2 n$ is tiny in comparison. This strongly suggests comparison to the series without this term.
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Geez, I should have seen that. My head was still spinning from the amount of algebra it took to get everything else out of the way. Thanks. – Trevor Alexander Feb 3 at 6:03
Hope it's not too far off topic but would like to follow up @T.Bongers' answer by sharing my favourite comparison of logarithms and powers. (And it's too long to fit in a comment.)
Consider the graph of $y=x^{1/10}$. You can calculate that if $x=1000000$ then $y$ is a bit less than $4$. Imagine drawing the graph on a scale of $1$ unit equals $1$ millimetre. Then although the graph is always increasing, it's like walking up a hill (if that's the right word) which climbs $4$ millimetres in a kilometre. This is clearly a climb which you would not even notice - you would think it was totally flat.
So, $y=x^{1/10}$ increases incredibly, almost unimaginably, slowly...
...but $y=\ln x$ increases even more slowly than that!! Certainly $\ln x$ is larger than $x^{1/10}$ for small values of $x$. If you draw both graphs on the same axis for say $1\le x\le1000$, you will get a very misleading impression - eventually, $x^{1/10}$ takes over and is larger than $\ln x$.
How long does it take for this to happen? Well, with a bit of trial and error you can find that the crossover point is approximately $x=4\times10^{15}$. If you (try to) draw this at a scale of $10$ centimetres for a unit - suitable for a blackboard in class - you find that the crossover point is about $3000$ times further away than the sun!
Hope this helps to illustrate how unimaginably slowly (I didn't say "almost" this time) the logarithm function grows.
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That's a good one. Closer to home for me is visual/auditory senses, which give you great sensitivity at low ranges without causing a seizure when a truck passes :) – Trevor Alexander Feb 3 at 23:20
Lemma $\displaystyle \sum_{n=2}^{\infty} \dfrac{\log_2 n}{n}$ diverges.
Proof By comparison to $\displaystyle \sum_{n=2}^{\infty} \dfrac{1}{n}$.
For $n\geq 2$, $\log_2 n \geq 1$. So $\dfrac{\log_2 n}{n} \geq \dfrac{1}{n}$. Since $\displaystyle \sum_{n=2}^{\infty} \dfrac{1}{n}$ diverges, $\displaystyle \sum_{n=2}^{\infty} \dfrac{\log_2 n}{n}$ diverges.
- | 2014-12-23T01:51:31 | {
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http://math.stackexchange.com/questions/157634/combinatorics-of-coin-tosses | # combinatorics of coin tosses
12 tosses of a fair coin result in an outcome of 6 heads and 6 tails.
How many sequences of such a result can have at most 4 successive outcomes of one kind ?
I have obtained what I believe is the correct answer as 924 - 84 = 840, by considering various sequences to be excluded
(a) 6-6 of a kind (b) 6-5 of a kind (c) 5-5 of a kind (d) only one 6 of a kind (e) only one 5 of a kind
However, this seems a very crude and tedious method. I am seeking help whether this can be done by inclusion-exclusion, or by some other slick way.
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Here is yet another way. It is still applicable when the number of H's and T's are changed or the maximal number $4$ in the question is changed:
Given a sequence of H's and T's, for example HTHHHHTTHTTT, note that H's and T's are partitioned. In the example, the (ordered) partition $p_H$ for H's is $1 + 4 + 1$, whereas it is $p_T = 1 + 2 + 3$ for T's. Note that the number of parts $l(p_H)$ and $l(p_T)$ of $p_H$ and $p_T$ are either the same or differ by 1. In the example, they are both 3. If they are the same, there are two possible ways to interlace H's and T's. If they differ by 1, there is only one way.
So, consider the partitions of 6. 6 = 6
= 5 + 1 = 4 + 2 = 3 + 3
= 4 + 1 + 1 = 3 + 2 + 1 = 2 + 2 + 2
= 3 + 1 + 1 + 1 = 2 + 2 + 1 + 1
= 2 + 1 + 1 + 1 + 1
= 1 + 1 + 1 + 1 + 1 + 1
Taking the order into account, there are 3 partitions having 2 parts, 10 partitions having 3 or 4 parts, 5 partitions having 5 parts and finally 1 partition having 6 parts, satisfying the criterion that no TTTTT and HHHHH must show up.
For $l(p_H), l(p_T)$ = 2, 2 we have $2.3.3 = 18$ sequences
For $l(p_H), l(p_T)$ = 2, 3 or 3, 2 we have $2.3.10 = 60$ sequences
For $l(p_H), l(p_T)$ = 3, 3 we have $2.10.10 = 200$ sequences
For $l(p_H), l(p_T)$ = 3, 4 or 4, 3 we have $2.10.10 = 200$ sequences
For $l(p_H), l(p_T)$ = 4, 4 we have $2.10.10 = 200$ sequences
For $l(p_H), l(p_T)$ = 4, 5 or 5, 4 we have $2.5.10 = 100$ sequences
For $l(p_H), l(p_T)$ = 5, 5 we have $2.5.5 = 50$ sequences
For $l(p_H), l(p_T)$ = 5, 6 or 6, 5 we have $2.5.1 = 10$ sequences
and finally
For $l(p_H), l(p_T)$ = 6, 6 we have $2.1.1 = 2$ sequences
So there are $18 + 60 + 200 + 200 + 200 + 100 + 50 + 10 + 2 = 840$ sequences.
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I really like the generalisation. I was unsuccessfully toying with the idea of using "stars and bars" somehow, which has materialised in an "interleaving" avatar, which can be used to work out inclusions or exclusions, as convenient. Thanking you and all others who have extended their valuable help. – true blue anil Jun 14 '12 at 13:00
Here's another solution. It's not very slick, but it works.
First, let's find the number of arrangements with 5 or more consecutive heads. Consider the eight symbols: T, 5T, H,H,H,H,H,H. These can be put into 8!/6!=56 distinct orders. For 14 of these, the symbols T and 5T are neighbors. Now, for example if we expand 5T into TTTTT in either "5T T HHHHHH" or "T 5T HHHHHH" we get the same result. Compensating for that, the number of arrangements with 5 or more consecutive heads is 42+14/2 = 49. Similarly, the number of arrangements with 5 or more consecutive tails is 49.
How many arrangements have both 5 or more consecutive heads and 5 or more consecutive tails? We can arrange the four symbols 5T,T,5H,H in 24 different ways.
1. For eight of these arrangements, the heads and tails are separated, e.g., T 5H 5T H. Each of these counts once when expanded.
2. For eight of these arrangements, either heads or tails are separated but not both, e.g., T 5H H 5T. Each of these counts twice when expanded.
3. For eight of these arrangements, neither the heads nor tails are
separated, e.g., H 5H 5T T. Each of these counts four times when expanded.
Therefore, the number of arrangements with both 5 or more consecutive heads and 5 or more consecutive tails is 8/1+8/2+8/4=14.
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Relooking at your approach, i find that it can be generalised for n heads, n tails with a max of (n-2) succesive outcomes of a kind as – true blue anil Jun 19 '12 at 6:59
(continued) C(2n,n) - 2(n+1)^2 + 14, which makes this approach extremely useful. Thanks ! – true blue anil Jun 19 '12 at 7:09
I'm glad that you find it useful. It's a fun problem! – Byron Schmuland Jun 19 '12 at 12:35
Here's another way. Let $a_n$ be the number of sequences in $n$ tosses of a coin in which there are at most 4 consecutive outcomes of one kind. Then by simple counting $a_1=2$, $a_2=4$, $a_3=8$, $a_4=16$, and $a_5=30$. With a little thought you can get the recurrence relation $$a_n=2a_{n-1}-a_{n-5}$$ Now there are two ways you can go. You can just use this recurrence to grind out $a_6=2\times30-2=58$, $a_7=2\times58-4=112$, and so on until you get up to $a_{12}$. Or you can use standard techniques to solve the (homogeneous, linear, constant-coefficient) recurrence, but I think it gets messy for this one. Well, that's what computers are for!
EDIT: OP has remembered that the outcome is to be 6 heads and 6 tails, so the above is irrelevant. So let's do inclusion-exclusion instead.
We start with $12\choose6$. We subtract all the ways of having 5 (or more) consecutive heads. That's 56, because the first of the 5 consecutive heads can be in any one of 8 locations, and, once you've fixed the first of the 5 consecutive heads, you've got 7 remaining tosses, of which 1 is supposed to be heads. There are also 56 ways of having (at least) 5 consecutive tails; subtract those as well.
Now add back in all the ways of having two instances of 5 consecutive heads. The only way for that to happen is to have 6 consecutive heads, and there are 7 ways for that to happen. Similarly, 7 ways to have two instances of 5 consecutive tails.
Also add in all ways of having 5 consecutive heads and 5 consecutive tails. If the heads are 1 through 5, there are 3 places for the tails, and 2 ways to fill in the other two slots; if the heads are 2 through 6, 2 places for tails, 2 ways to fill; heads 3 through 7, 1 place for tails, 2 ways to fill, making 12 ways with the runs of heads preceding the run of tails; another 12 the other way around.
Now we have to subtract out all the ways that meet three conditions. By the way, if you're getting the impression that there's nothing slick about inclusion-exclusion, at least for this problem, you're right, but it's what the customer wanted, and we've come this far, so let's go for it. We can have 6 consecutive heads and (at least) 5 consecutive tails, 4 ways (heads starting at 1, 2, 6, or 7), or 6 consecutive tails and 5 (or more) consecutive heads, another 4 ways.
Finally, we have to add back in the 2 ways of meeting four conditions; 6 heads followed by 6 tails, or the other way around.
So the answer would appear to be $${12\choose6}-56-56+7+7+12+12-4-4+2$$ whatever that is.
MORE EDIT: I figured out where I went wrong. In the next-to-last case, meeting three conditions, 6 heads and 5 (or more) tails can happen 6 ways, not 4. With the 6 heads starting at 1, the 5 tails can start at 7 or 8, and with the six heads starting at 7, the five tails can start at 1 or 2. Add in the heads starting at 2 or at 6 and you get 6 ways, not 4. This makes the final calculation $${12\choose6}-56-56+7+7+12+12-6-6+2$$ which is 840, in agreement with everyone else.
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The procedure you suggest gives an answer of 2980, whereas all possible sequences are only C(12,6) = 924. There seems to be a communication gap somewhere ! – true blue anil Jun 13 '12 at 7:50
@Gerry Myerson : I'm so sorry. In the original question, I needed to mention that the outcome of the 12 tosses is 6 heads and 6 tails. – true blue anil Jun 13 '12 at 8:50
So, what are you waiting for? Edit that into the question! Now! – Gerry Myerson Jun 13 '12 at 10:08
@Gerry: I got a count of 840 against 844 that you are getting. Let me recheck. The answers should tally. – true blue anil Jun 13 '12 at 11:37
@Gerry: I have rechecked and get the same # for exclusions as before. (a) 6-6 of a kind:2 (b) 6-5 of a kind: 4 (c) 5-5 of a kind:8 (d) only one 6 of a kind: 6 (e) only one 5 of a kind: 64. 924 - 84- 840.The only count that is tedious is the last one. I had written "inclusion-exclusion or by some other slick way", hoping that I-E would be slick. I wonder if there is a slick way ? – true blue anil Jun 13 '12 at 12:22 | 2014-10-21T07:38:48 | {
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https://math.stackexchange.com/questions/2810715/spivaks-calculus-proofs-concerning-pascals-triangle | # Spivak's Calculus: Proofs concerning Pascal's Triangle
Problem 3 of Chapter 2 in Spivak's Calculus poses 5 problems associated with Pascal's Triangle. The first of these asks you to prove that $\binom{n+1}{k}$ = $\binom{n}{k-1} + \binom{n}{k}$ which was fairly straightforward.
The second task was to prove by induction that $\binom{n}{k}$ is always a natural number. For this, I took the recursion approach: given the demonstration of the equality above, any $\binom{n}{k}$ can be decomposed into $\binom{n-1}{k-1} + \binom{n-1}{k}$ and so on, with each value decomposing until every binomial reaches either an $n$ or $k$ value of $0$, at which point they evaluate to 1 and can be summed. $1$ is a natural number, so the sum of any number of $1$s must therefore be a natural number. This proof seems to me as though it lacks formality, but I think the logic is appropriate.
The third task, however, I am simply unsure what to make of. It reads: "Give another proof that $\binom{n}{k}$ is a natural number by showing that $\binom{n}{k}$ is the number of sets of exactly $k$ integers each chosen from $1, \cdots,n$." The word induction is not used in the prompt, and I can see that this appears to be true across the first few rows of the triangle, but how one would rigorously prove something like this, I do not know. Additionally, having proven that this is true, is the link from that to $\binom{n}{k}$ always being a natural number just the fact that you cannot have partial or negative sets, and therefore it must be a natural number?
• I think the exercise want you to prove that $c(n,k)$ the number of sets of exactly $k$ integers chosen among $n$ follows the same induction formula, i.e. $c(n,k)=c(n-1,k-1)+c(n-1,k)$ and since you already noticed it takes the same values in the first rows of the triangle then you deduce $c(n,k)=\binom{n}{k}$. – zwim Jun 6 '18 at 22:49
By definition ${n\choose k}=\frac{n!}{k!(n-k)!}=\frac{n\cdot (n-1)\cdot\ldots\cdot (n-k+1)}{k!}=\frac{n}{k}\cdot\frac{n-1}{k-1}\cdot\ldots\cdot\frac{n-k+1}{1}$.
Now consider $\{1,\ldots,n\}$. There are $n$ possibilities to pick an integer from this set. Let this integer be called $a_1$.
Now there are $n-1$ possibilities to pick an element from the set $\{1,\ldots,n\}\setminus \{a_1\}$.
By iterating this procedure we obtain $n\cdot (n-1)\cdot\ldots\cdot (n-k+1)$ possibilities to choose $(a_1,\ldots,a_k)$ from $\{1,\ldots,n\}$. Now note that we counted some possibilities multiple times, for example $a_1=1,a_2=2$ and $a_1=2,a_2=1$ result to the same integers chosen. So we have to divide by the number of possibilities to order $k$ elements which is $k!$ (there are $k$ possibilities to place the first integer, $k-1$ to place the second etc.).
• Very helpful, thank you. Is there some any special study I should undertake to develop the ability to notice this, or does it fall under the general heading of 'mathematical maturity' which one acquires through sheer experience? – user242007 Jun 6 '18 at 23:11
• Well I think the best way to solve these kind of problems is to first work out the combinatorical side (here the number of possibilities to choose $k$ elements from $\{1,\ldots,n\}$) and then try to rearrange the left side to find this number. – Nightgap Jun 6 '18 at 23:26
• To the proposer: The topic of this Answer is Finite Combinatorics. There is a nice book called Finite Mathematics by Kemeny, Snell, & Thompson with many different topics (including combinatorics) utilizing a variety of techniques. Example: An election with 2 candidates A & B..... "A" won by $m$ votes to $n.$ When the ballots are counted, one at a time, what are the odds that A is always ahead during the count? – DanielWainfleet Jun 7 '18 at 0:35
Was writing out something along the lines of what's already so nicely explained by @Nightgap :), but just for good measure and I hope you don't mind that it's not really an answer to your specific question posed above, will still post what was going to be an addendum which has a different flavour and is nice to know if you haven't seen it (came across it in Alan Baker's Concise Intro to Number Theory aeons ago - great little book!).
So let $\lfloor x\rfloor$ be the largest integer $\leq x$. Let $p$ be a prime, then there $\lfloor n/p \rfloor$ elements of the set $\{1,2,...n\}$ divisible by $p$, $\lfloor n/p^2 \rfloor$ divisible by $p^2$, and so on. This means that the highest power $k$ of $p$ such that $p^k |n!$ is given by:$$k=\sum_{j=1}^{\infty}\lfloor n/p^j\rfloor \qquad \qquad (1)$$
Essentially you're adding to the sum each time you go up a prime power, so eg first you pick up all the unit powers with the term $\lfloor n/p \rfloor$, but then you need to add to these the extra contribution of the prime square multiples with the $\lfloor n/p^2 \rfloor$ term, then for the cubes, and so on. So now we have an explicit expression for the maximal power of a prime $p^k$ such that it divides $n!$
Now note that in general: $$\lfloor x+y \rfloor \geq \lfloor x\rfloor + \lfloor y \rfloor$$ (eg if the fractional parts of $x$ and $y$ add up to an integer), hence: $$\lfloor n/p^j \rfloor \geq \lfloor (n-k)/p^j \rfloor + \lfloor k/p^j \rfloor \qquad \qquad (2)$$
But then: $${n \choose k}=\frac{n!}{k!(n-k)!}$$ ...must be a natural number since by (1) and (2) all the prime powers in the denominator will be 'absorbed' by those in the numerator. Hope that makes sense and again though not how you would answer your exercise I think is a nice perspective to have on binomial coefficients, in addition to the all important combinatorial one. Cheers.
• An interesting alternative look at the problem. It's going to take some time to digest this, but I greatly appreciate you taking the time to offer another angle. – user242007 Jun 7 '18 at 0:18
• No problem! Like I say you can find it in A Concise Introduction to Number Theory by Alan Baker, well worth checking out if you have the time, but for me the above was a pretty powerful aid to understanding the structure of binomial coefficients from a number theoretic perspective. – Mehness Jun 7 '18 at 0:20
• Just to add the above proof also immediately shows that multinomial coefficients $\frac{n!}{k_1! \cdot k_2! \cdots k_r!}$ where $\sum_{i=1}^r k_i = n$ are also natural numbers – Mehness Jun 7 '18 at 10:31
For convenience, let $S_n = \{x_1, x_2, x_3, \dots, x_n\}$ represent a set of $n$ distinct objects.
Lets use $_nC_k$ to represent the number of ways of selecting $k$ objects from $S_n$.
So what does $_{n+1}C_k$ look like?
An element of such a selection will either include the object $x_{n+1}$ or it will not. We already know that there are $_nC_k$ ways to choose $k$ objects from the objects $x_1, x_2, \dots, x_n$. So how do we count the number of ways to select $k$ objects that contain $x_{n+1}$. Easy; we select $k-1$ objects from $S_n$ {there are $_nC_{k-1}$ such ways} and then add $x_{n+1}$ to the collection.
It follows that $_nC_{k-1} + _nC_k = _{n+1}C_k$. | 2019-06-18T20:39:40 | {
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https://math.stackexchange.com/questions/342128/for-every-integer-n-15-mid-n-iff-3-mid-n-and-5-mid-n | # For every integer $n$, $15\mid n$ iff $3\mid n$ and $5\mid n$
I'm trying to prove that for every integer $n$, $15\mid n$ iff $3\mid n$ and $5\mid n$. The first part of this bi-conditional was easy for me to prove, but I'm having problems with the second. Here is what I have so far:
Suppose $15\mid n$. Then we can let $k$ be some integer such that $15k=n$. Since $3(5k)=15k=n, 3\mid n$. Also, since $5(3k)=15k=n$, it follows that $5\mid n$.
Now I'm stuck trying to prove it the other way. I've assumed that for some integers $k$ and $j$, $3k=n$ and $5j=n$, and am trying to come up with something of the form $15p=n$ for some integer $p$, but I can only seem to come up with fractions, which seems wrong. Any hints?
It is a good start, not much more is needed.
You have $3k=n$ and $5j=n$. It follows that $3k=5j$.
Note that $3k=5j$ in particular says that $3$ divides $5j$. It is a standard fact of number theory that if a prime $p$ (here $p=3$) divides a product $ab$, then $p$ divides $a$ or $p$ divides $b$ (or both). Clearly $3$ does not divide $5$. So $3$ divides $j$.
Let $j=3w$. Then $n=5j=(5)(3w)=15w$.
Remark: $1.$ The fact that if a prime $p$ divides $ab$, then $p$ divides $a$ or $p$ divides $b$ is sometimes called Euclid's Lemma. Books VII to IX, and to a fair degree Book X of Euclid's Elements are number-theoretic in character.
$2.$ We do not really need Euclid's Lemma in this case. For we only want to show that if $3$ divides $5j$, then $3$ divides $j$. If $3$ does not divide $j$, then $j$ leaves a remainder of $1$ or $2$ on division by $3$. In either case, you can show that the remainder when $5j$ is divided by $3$ is not $0$.
• Ahh, Thanks! This problem is out of How to Prove it by Velleman, and I haven't actually had a chance to start studying number theory. I'll definitely remember this, thanks! – stochasm Mar 26 '13 at 22:07
• You are welcome. I added in the remarks a little something about doing it without euclid's Lemma, or more precisely proving euclid's Lemma in the case $p=3$. So we can actually write a detailed proof with no reference to other stuff. – André Nicolas Mar 26 '13 at 22:20
Hint $\rm\ \ 5,3\mid n\:\Rightarrow\: \dfrac{n}5,\,\dfrac{n}3\in\Bbb Z\ \Rightarrow\ 2\left(\dfrac{n}5\right)-\dfrac{n}3 \, =\, \dfrac{n}{15}\in\Bbb Z\:\Rightarrow\: 15\mid n$
$d \mid n$ means $\exists k, dk = n$.
So we want to prove that $\exists k, 15k = n$ iff $\exists h,l, 5h=3l=n$.
To do one direction write $15 = 3 \cdot 5$: there $\exists k, 3 \cdot 5 \cdot k = n$ so $5(3k)=3(5k)=n$ proves the implication with $h=3k$, $l=5k$.
In the other direction, we have $5h=3l=n$ so $3\mid h$ (and $5 \mid l$), put $k = h/3$ and you get $5 \cdot 3 \cdot k = n$.
Since $3|n$ then $n=3k$. But you also know $5|n$, that is, $5|3k$. Now, here use the fact that $5\not|3$.
By using that fact you get that $5|k$, so $k=5j$. Therefore $n=3k=3(5j)=15j$.
Now $15$ divides $3n$ and $5n$, and $n = 10n - 9n = 2 \cdot (5n) - 3 \cdot (3n)$ so..
This is basically using Bezout and the fact that $3$ and $5$ are coprime. When $a$ and $b$ are coprime, $xa + yb = 1$ for some integers $x$ and $y$. From this a more general theorem can be proven.
If $a$ and $b$ are coprime integers, both dividing $n$, then $n$ is divisible by $ab$.
Proof: $an$ and $bn$ are divisible by $ab$, so $n = xan + ybn$ is divisible by $ab$.
Your problem is the case $a = 3$, $b = 5$. | 2019-11-15T21:12:50 | {
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https://math.stackexchange.com/questions/2205612/surface-area-of-a-sphere-with-integration-of-disks | # Surface area of a sphere with integration of disks [duplicate]
Why it is not correct to say that the surface area of a sphere is:
$$2 \int_{0}^{R} 2\pi r \text{ } dr$$
In my mind we are summing up the perimeters of disks from $r=0$ to $r=R$, so by 1 integration, we would have $\frac{1}{2}$ of the surface area of the sphere.
I know that it's not correct because that will give us $2\pi R^2$ that it's different from $4\pi R^2$, but why???
Thanks!
• It's bad form, and technically incorrect, to have one of the limits of integration depend on the variable of integration. What does it mean to say "to $r=r$"? Isn't $r=r$ always true since $r$ varies between $0$ and some fixed positive constant? It'd be better to call this constant $R$ to keep it separate from the variable $r$. This is more of a side note. – user307169 Mar 27 '17 at 17:37
• @tilper fixed It, thanks – Bruno Reis Mar 27 '17 at 17:41
• If your method was correct, then you should get the same answer for a cone whose base radius and height are equal. But you don't, do you? – uniquesolution Mar 27 '17 at 17:54
In order to give a short explanation of the mistake :
In writing $\quad 2 \int_{0}^{R} 2\pi r \text{ } dr \quad$ you sum elementary areas $\quad dA=2\pi\:r\:dr\quad$
that is to say, a strip length$=2\pi r$ and width$=dr$.
The hitch is that the width is not $dr$ but is $ds$ (see the figure).
With a few calculus, you can find $\quad ds=\frac{R}{\sqrt{R^2-r^2}}dr \quad\to\quad dA=2\pi\:r\frac{R}{\sqrt{R^2-r^2}}dr$
$$\quad 2 \int_{0}^{R} 2\pi\:r\frac{R}{\sqrt{R^2-r^2}}dr =4\pi R^2$$
This is a method to find the area of the sphere. Of course, a simpler method consists in doing the job in spherical coordinates instead of Cartesian coordinates.
• Perfect explanation! Thanks! – Bruno Reis Mar 27 '17 at 23:22
What you are doing is indeed integrating over rings, but then you get the area of a full circle, which was the result of your calculation up to the irrelevant factor of 2.
You can imagine that by adding more and more dense rings you fill a circle with radius $R$ (the upper limit of your integral).
But to get the area of a sphere, you need to parametrize all 3 dimensions. What you can do is use spherical coordinates, where each point in space is described by a radius $r$, and two angles $\theta,\phi$ as in the link.
In fact, to get the area of the sphere you need to keep the radius constant, and integrate over the angles that parametrize the given sphere.
The result yields
$$S=\int_0^{\pi}\int_0^{2\pi}R^2\sin\theta d\phi d\theta=R^22\pi\int_0^{\pi}\sin\theta d\theta=4\pi R^2$$ as it should be.
A) when you pile up the disks you shall pay attention to law with which $r$ varies (with $h$, $x$, $\alpha$ $\cdots$) otherwise, as already told you do not make any distinction with a cone, a cylinder,..
B) adding disks to find volumes is ok, because the volume of the triangular cross-section ring that is left-out is a higher order infinitesimal vs. the volume of the disk; its surface instead is of the same order, and you must consider the hypothenuse instead of the cathetus.
$$\text{Area will be}~~: 2 \int_{0}^{R} 2\pi x ~ ds$$
Where $ds$ is width of strip bounded by circles of radius $x$ and $x+dx$ situated at height $y$. Also $ds \neq dr$ it's tilted in $y$ direction too. Only horizontal projection of $ds$ is $dr$.What you have done is valid for Disk See image below.
$$(ds)^2=(dx)^2+(dy)^2$$
Therefore : $$ds = \sqrt{1+\Bigg( \frac{dy}{dx}\Bigg)^2}$$ Now,
$$x^2 + y^2 = R^2$$ $$y = \sqrt{R^2 - x^2}$$ $$\dfrac{dy}{dx} = \dfrac{-2x}{2\sqrt{R^2 - x^2}}$$ $$\dfrac{dy}{dx} = \dfrac{-x}{\sqrt{R^2 - x^2}}$$ $$\left( \dfrac{dy}{dx} \right)^2 = \dfrac{x^2}{R^2 - x^2}$$
Thus, $$\displaystyle A = 4\pi \int_0^R x \sqrt{1 + \dfrac{x^2}{R^2 - x^2}} \, dx$$ $$\displaystyle A = 4\pi \int_0^R x \sqrt{\dfrac{(R^2 - x^2) + x^2}{R^2 - x^2}} \, dx$$ $$\displaystyle A = 4\pi \int_0^R x \sqrt{\dfrac{R^2}{R^2 - x^2}} \, dx$$
Let $$x = R \sin θ \implies dx = R \cos θ dθ$$
When $x = 0, θ = 0$ When $x = R, θ = \pi/2$
Thus, $$\displaystyle A = 4\pi \int_0^{\pi/2} R \sin \theta \sqrt{\dfrac{R^2}{R^2 - R^2 \sin^2 \theta}} \, (R \cos \theta \, d\theta)$$ $$\displaystyle A = 4\pi \int_0^{\pi/2} R^2 \sin \theta \cos \theta\sqrt{\dfrac{R^2}{R^2(1 - \sin^2 \theta)}} \, d\theta$$ $$\displaystyle A = 4\pi R^2 \int_0^{\pi/2} \sin \theta \cos \theta\sqrt{\dfrac{1}{\cos^2 \theta}} \, d\theta$$ $$\displaystyle A = 4\pi R^2 \int_0^{\pi/2} \sin \theta \cos \theta \left( \dfrac{1}{\cos \theta} \right) \, d\theta$$ $$\displaystyle A = 4\pi R^2 \int_0^{\pi/2} \sin \theta \, d\theta$$ $$A = 4\pi R^2 \bigg[-\cos \theta \bigg]_0^{\pi/2}$$ $$A = 4\pi R^2 \bigg[-\cos \frac{1}{2}\pi + \cos 0 \bigg]$$ $$A = 4\pi R^2 \bigg[ -0 + 1 \bigg]$$ $$A = 4\pi R^2$$
If you consider a formula like
$$dA = 2 \pi r \, dr$$
for a cylinder, it is zero and so fails.
If you consider a formula for a cylinder,
$$dA = 2 \pi r \, dz$$
it seems to be correct.
However, for a cone $dz=0$ so even this formula fails.
Intuition says $dr,dz$ should both be involved.
In fact if you recall cone formula $A = \pi\,r\, s = 2 \pi\, \bar {r} \, s ,\, dA= 2\pi r ds$
slant generator is correct in principle.
We have cone differential area
$$dA = 2 \pi r ds = 2 \pi r \sqrt{dr^2+dz^2 }$$
and that should be considered in integration of all surfaces of revolution.
There are two problems here.
First: Not only do your disks change radius from the apex to the base of the hemisphere, but they do so at a varying rate. Let's say (for concreteness) that the sphere is $$x^2 + y^2 + z^2 = R^2$$, and we're dividing the sphere into disks based on the $$x$$-coordinate. If you try to approximate the surface area of a hemisphere of radius $$R=1$$ by slicing it into disks of radius (say) $$0.1, 0.2, \ldots, 0.9, 1$$, you'll find that the first few small-radius slices are very close together— to be precise, at $$x = \sqrt{1 - 0.1^2} = 0.994$$, $$\sqrt{1-0.2^2} = 0.980$$, $$\sqrt{1-0.3^2} = 0.954$$, and so forth—but the larger slices are much farther apart. If you use the surface areas of these disks to calculate the surface area of the sphere, you have to take into account the fact that the disks have different widths.
Going from finite to infinitesimal width: The radius of a disk at $$x$$ is $$r = \sqrt{R^2-x^2}$$ (the inverse of this formula is $$x = \sqrt{R^2 - r^2}$$), so if you make cuts in the sphere each time the radius changes by $$dr$$, then the slices will have width $$-dx = -dr \dfrac{dx}{dr} = \dfrac{r}{\sqrt{R^2-r^2}}\, dr$$ (the negative sign corrects for the fact that $$x$$ increases with decreasing $$r$$).
Second: The surfaces of the sphere slices that we discuss above have sloping cross-sections (in the $$xy$$ plane or any plane that includes the $$x$$-axis), whereas cylinders have flat surface cross-sections. You can ignore the difference this makes for calculating volume, but not for surface area; otherwise you can prove that $$\pi = 4$$. So you need to make another correction: the width of the surface of a disk isn't $$dx$$ but rather $$\sqrt{(dx)^2 + (dr)^2} = \sqrt{ \dfrac{r^2}{R^2 - r^2} + 1}\, dr = \dfrac{R}{\sqrt{R^2 - r^2}}\, dr$$. Putting this all together and integrating over $$r$$ (though $$x$$ would be nicer): The integral that you actually have to do is thus \begin{align*} &\int_0^R (2\pi r) \frac{R}{\sqrt{R^2 - r^2}}\, dr \\ =& \;2\pi R \int_{r=0}^{r=R} \frac{-du/2}{\sqrt{u}} \tag{u = R^2 - r^2; du = -2r\, dr} \\ =& -2\pi R \left. \sqrt{u} \right|_{r=0}^{r=R} \\ =& -2\pi R \left. \sqrt{R^2 - r^2}\right|_{r=0}^{r=R} \\ =&\;2\pi R^2 \end{align*} and this is the result you want. | 2021-02-26T02:26:49 | {
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https://vioora.com/hugh-frater-qochpg/f12ce6-multiplying-complex-numbers-in-rectangular-form | # multiplying complex numbers in rectangular form
(5 + j2) + (2 - j7) = (5 + 2) + j(2 - 7) = 7 - j5 (2 + j4) - (5 + j2) = (2 - 5) + j(4 - 2) = -3 + j2; Multiplying is slightly harder than addition or subtraction. Ask Question Asked 1 year, 6 months ago. Viewed 385 times 0 $\begingroup$ I have attempted this complex number below. Addition of Complex Numbers . Here are some specific examples. Find the complex conjugate of the denominator, also called the z-bar, by reversing the sign of the imaginary number, or i, in the denominator. As imaginary unit use i or j (in electrical engineering), which satisfies basic equation i 2 = −1 or j 2 = −1.The calculator also converts a complex number into angle notation (phasor notation), exponential, or polar coordinates (magnitude and angle). Now, let’s multiply two complex numbers. The standard form, a+bi, is also called the rectangular form of a complex number. In other words, to write a complex number in rectangular form means to express the number as a+bi (where a is the real part of the complex number and bi is the imaginary part of the complex number). So just remember when you're multiplying complex numbers in trig form, multiply the moduli, and add the arguments. How to Write the Given Complex Number in Rectangular Form". The major difference is that we work with the real and imaginary parts separately. The rectangular form of a complex number is written as a+bi where a and b are both real numbers. Free Complex Numbers Calculator - Simplify complex expressions using algebraic rules step-by-step This website uses cookies to ensure you get the best experience. if you need any other stuff in math, please use our google custom search here. Free Complex Number Calculator for division, multiplication, Addition, and Subtraction. Complex numbers in the form are plotted in the complex plane similar to the way rectangular coordinates are plotted in the rectangular plane. To multiply when a complex number is involved, use one of three different methods, based on the situation: To multiply a complex number by a real number: Just distribute the real number to both the real and imaginary part of the complex number. Plot each point in the complex plane. and x − yj is the conjugate of x + yj.. Notice that when we multiply conjugates, our final answer is real only (it does not contain any imaginary terms.. We use the idea of conjugate when dividing complex numbers. This calculator does basic arithmetic on complex numbers and evaluates expressions in the set of complex numbers. Multipling and dividing complex numbers in rectangular form was covered in topic 36. Complex numbers can be expressed in numerous forms. Label the x-axis as the real axis and the y-axis as the imaginary axis. Hence the Re (1/z) is (x/(x2 + y2)) - i (y/(x2 + y2)). However, due to having two terms, multiplying 2 complex numbers together in rectangular form is a bit more tricky: (3z + 4zbar â 4i) = [3(x + iy) + 4(x + iy) bar - 4i]. Yes, you guessed it, that is why (a+bi) is also called the rectangular form of a complex number. Recall that FOIL is an acronym for multiplying First, Outer, Inner, and Last terms together. A1. 1. Multiplication . 1. In essence, the angled vector is taken to be the hypotenuse of a right triangle, described by the lengths of the adjacent and opposite sides. Example 7 MULTIPLYING COMPLEX NUMBERS (cont.) It is the distance from the origin to the point: See and . Simplify. To understand and fully take advantage of multiplying complex numbers, or dividing, we should be able to convert from rectangular to trigonometric form and from trigonometric to rectangular form. Math Gifs; Algebra; Geometry; Trigonometry; Calculus; Teacher Tools; Learn to Code; Home; Algebra ; Complex Numbers; Complex number Calc; Complex Number Calculator. Note that the only difference between the two binomials is the sign. Find powers of complex numbers in polar form. Real numbers can be considered a subset of the complex numbers that have the form a + 0i. The different forms of complex numbers like the rectangular form and polar form, and ways to convert them to each other were also taught. So 18 times negative root 2 over. In polar form, the multiplying and dividing of complex numbers is made easier once the formulae have been developed. This is an advantage of using the polar form. Find the complex conjugate of the denominator, also called the z-bar, by reversing the sign of the imaginary number, or i, in the denominator. 7) i 8) i Use rectangular coordinates when the number is given in rectangular form and polar coordinates when polar form is used. Polar form is where a complex number is denoted by the length (otherwise known as the magnitude, absolute value, or modulus) and the angle of its vector (usually denoted by … Included in the resource: 24 Task cards with practice on absolute value, converting between rectangular and polar form, multiplying and dividing complex numbers … Label the x-axis as the real axis and the y-axis as the imaginary axis. Example 4: Multiplying a Complex Number by a Real Number . To add complex numbers in rectangular form, add the real components and add the imaginary components. Let z 1 = r 1 cis θ 1 and z 2 = r 2 cis θ 2 be any two complex numbers. Using either the distributive property or the FOIL method, we get A complex number can be expressed in standard form by writing it as a+bi. $\text{Complex Conjugate Examples}$ $\$$3 \red + 2i)(3 \red - 2i) \\(5 \red + 12i)(5 \red - 12i) \\(7 \red + 33i)(5 \red - 33i) \\(99 \red + i)(99 \red - i) This is the currently selected item. Although the complex numbers (4) and (3) are equivalent, (3) is not in standard form since the imaginary term is written first (i.e. In the complex number a + bi, a is called the real part and b is called the imaginary part. Complex numbers in the form are plotted in the complex plane similar to the way rectangular coordinates are plotted in the rectangular plane. Multiplication of Complex Numbers. To add complex numbers, add their real parts and add their imaginary parts. Rectangular Form of a Complex Number. Rectangular Form of a Complex Number. This material appears in section 6.5. In other words, to write a complex number in rectangular form means to express the number as a+bi (where a is the real part of the complex number and bi is the imaginary part of the complex number). The primary reason for having two methods of notation is for ease of longhand calculation, rectangular form lending itself to addition and subtraction, and polar form lending itself to multiplication and division. This is an advantage of using the polar form. Change ), You are commenting using your Google account. Post was not sent - check your email addresses! The following development uses trig.formulae you will meet in Topic 43. In other words, there are two ways to describe a complex number written in the form a+bi: To write a complex number in rectangular form you just put it into the standard form of a complex number by writing it as a+bi. How to Divide Complex Numbers in Rectangular Form ? There are two basic forms of complex number notation: polar and rectangular. When in rectangular form, the real and imaginary parts of the complex number are co-ordinates on the complex plane, and the way you plot them gives rise to the term “Rectangular Form”. 2 and 18 will cancel leaving a 9. B2 ( a + bi) Error: Incorrect input. Addition, subtraction, multiplication and division can be carried out on complex numbers in either rectangular form or polar form. How do you write a complex number in rectangular form? Hence the value of Im(3z + 4zbar â 4i) is - y - 4. Here we are multiplying two complex numbers in exponential form. bi+a instead of a+bi). Multiplying by the conjugate . Subtraction is similar. Divide complex numbers in rectangular form. It will perform addition, subtraction, multiplication, division, raising to power, and also will find the polar form, conjugate, modulus and inverse of the complex number. The symbol ' + ' is treated as vector addition. https://www.khanacademy.org/.../v/polar-form-complex-number That’s right – it kinda looks like the the Cartesian plane which you have previously used to plot (x, y) points and functions before. If z = x + iy , find the following in rectangular form. Rectangular Form A complex number is written in rectangular form where and are real numbers and is the imaginary unit. Multiplication of Complex Numbers. This can be a helpful reminder that if you know how to plot (x, y) points on the Cartesian Plane, then you know how to plot (a, b) points on the Complex Plane. 5) i Real Imaginary 6) (cos isin ) Convert numbers in rectangular form to polar form and polar form to rectangular form. In polar form, the multiplying and dividing of complex numbers is made easier once the formulae have been developed. Find (3e 4j)(2e 1.7j), where j=sqrt(-1). Answer. Note that all the complex number expressions are equivalent since they can all ultimately be reduced to -6 + 2i by adding the real and imaginary terms together. Let’s begin by multiplying a complex number by a real number. Multiplying a complex number by a real number is simple enough, just distribute the real number to both the real and imaginary parts of the complex number. You may have also noticed that the complex plane looks very similar to another plane which you have used before. polar form Introduction When two complex numbers are given in polar form it is particularly simple to multiply and divide them. To convert from polar form to rectangular form, first evaluate the trigonometric functions. In this lesson you will investigate the multiplication of two complex numbers v and w using a combination of algebra and geometry. Multiplying Complex Numbers Together. This video shows how to multiply complex number in trigonometric form. 18 times root 2 over 2 again the 18, and 2 cancel leaving a 9. Now that we can convert complex numbers to polar form we will learn how to perform operations on complex numbers in polar form. ; The absolute value of a complex number is the same as its magnitude. Products and Quotients of Complex Numbers; Graphical explanation of multiplying and dividing complex numbers; 7. Rectangular form, on the other hand, is where a complex number is denoted by its respective horizontal and vertical components. It is the distance from the origin to the point: See and . Rather than describing a vector’s length and direction by denoting magnitude and … How to Write the Given Complex Number in Rectangular Form : Here we are going to see some example problems to understand writing the given complex number in rectangular form. Fill in your details below or click an icon to log in: You are commenting using your WordPress.com account. What you can do, instead, is to convert your complex number in POLAR form: #z=r angle theta# where #r# is the modulus and #theta# is the argument. This screen shows how the TI–83/84 Plus displays the results found in parts (a), (b), and (d) in this example. That is, [ (a + ib)/(c + id) ] ⋅ [ (c - id) / (c - id) ] = [ (a + ib) (c - id) / (c + id) (c - id) ] Examples of Dividing Complex Numbers We sketch a vector with initial point 0,0 and terminal point P x,y . A complex number is a number of the form a + bi, where a and b are real numbers, and i is an indeterminate satisfying i 2 = −1.For example, 2 + 3i is a complex number. Draw a line segment from \(0$$ to $$z$$. d) Write a rule for multiplying complex numbers. Rectangular Form. z 1 z 2 = r 1 cis θ 1 . Multiplying a complex number by a real number is simple enough, just distribute the real number to both the real and imaginary parts of the complex number. We start with an example using exponential form, and then generalise it for polar and rectangular forms. The Complex Hub aims to make learning about complex numbers easy and fun. However, due to having two terms, multiplying 2 complex numbers together in rectangular form is a bit more tricky: Find quotients of complex numbers in polar form. Complex Number Lesson . Multiplication and division of complex numbers in polar form. Example 1. Therefore the correct answer is (4) with a=7, and b=4. When you’re dividing complex numbers, or numbers written in the form z = a plus b times i, write the 2 complex numbers as a fraction. For Example, we know that equation x 2 + 1 = 0 has no solution, with number i, we can define the number as the solution of the equation. Multiplying both numerator and denominator by the conjugate of of denominator, we get ... "How to Write the Given Complex Number in Rectangular Form". Complex numbers are numbers of the rectangular form a + bi, where a and b are real numbers and i = √(-1). I get -9 root 2. Notice the rectangle that is formed between the two axes and the move across and then up? Polar Form of Complex Numbers; Convert polar to rectangular using hand-held calculator; Polar to Rectangular Online Calculator ; 5. 2.3.2 Geometric multiplication for complex numbers. To write a complex number in rectangular form you just put it into the standard form of a complex number by writing it as a+bi. Converting a complex number from polar form to rectangular form is a matter of evaluating what is given and using the distributive property. Rectangular form, on the other hand, is where a complex number is denoted by its respective horizontal and vertical components. But then why are there two terms for the form a+bi? But complex numbers, just like vectors, can also be expressed in polar coordinate form, r ∠ θ . ( Log Out / To divide the complex number which is in the form (a + ib)/ (c + id) we have to multiply both numerator … It is no different to multiplying whenever indices are involved. To write complex numbers in polar form, we use the formulas and Then, See and . To multiply complex numbers in polar form, multiply the magnitudes and add the angles. ; The absolute value of a complex number is the same as its magnitude. Multiplication and division in polar form Introduction When two complex numbers are given in polar form it is particularly simple to multiply and divide them. B1 ( a + bi) A2. Complex Number Functions in Excel. Apart from the stuff given in this section "How to Write the Given Complex Number in Rectangular Form", if you need any other stuff in math, please use our google custom search here. To plot a complex number a+bi on the complex plane: For example, to plot 2 + i we first note that the complex number is in rectangular (a+bi) form. A complex number is a number of the form a + bi, where a and b are real numbers, and i is an indeterminate satisfying i 2 = −1.For example, 2 + 3i is a complex number. So the root of negative number √-n can be solved as √-1 * n = √ n i, where n is a positive real number. We know that i lies on the unit circle. Change ), You are commenting using your Twitter account. Show Instructions. Example 2(f) is a special case. A complex number can be represented in the form a + bi, where a and b are real numbers and i denotes the imaginary unit. Multiplying Complex Numbers. (This is because it is a lot easier than using rectangular form.) 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Form is as simple as multiplying and dividing of complex numbers imaginary axis be carried Out complex. Property or the FOIL method + 4zbar â 4i ) is - y - 4 yi in set., just like vectors, we first need some kind of standard mathematical notation \! Form or polar form, multiply the magnitudes and … Plot each in! With steps shown 2j is the rectangular form of complex number from form. Binomials is the sign does basic arithmetic on complex numbers, just like,. The distributive property made easier once the formulae have been developed can to. Expressed in polar form. 2.4 of the text for an introduction to complex,. Rectangular forms Euler Identity interactive graph ; 6 form Step 1 sketch a with! Binomials is the conjugate of x + yi in the form are plotted in the form plotted. Numbers to polar form to rectangular using hand-held calculator ; 5 is written in rectangular form of complex without! 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And Euler Identity interactive graph ; 6 multiply the moduli, and the! 2 units along the horizontal axis, followed by 1 unit up on the other hand, where. That i lies on the unit circle, multiply the moduli, and b=4 exponential forms we first need kind! The distance from the stuff given in rectangular form '' is where complex! Very similar to the way rectangular coordinates are plotted in the form a complex number.! Coordinates when polar form is used this section, we use the formulas and then generalise it for polar rectangular... They 're in polar form. 0\ ) to \ ( 0\ ) \... Sketch a vector with initial point 0,0 and terminal point P x, y imaginary parts be represented a! Search here two angles the conjugate of 3 − 2j given in rectangular form. on! Then up ) Error: Incorrect input Incorrect input ) is - y - 4 forms. Of all three four digit numbers formed with non zero digits in general: x iy! Form used to Plot complex numbers is easy in rectangular form and polar coordinates when polar to... A+Bi, is also called the imaginary axis because it is a special.. Easy and fun followed by 1 unit up on the complex plane similar to the way coordinates! ( 3e 4j ) ( 2e 1.7j ), you are commenting using your Twitter account angle θ ” )! These complex numbers sum of all three four digit numbers formed with non zero digits root... R 2 cis θ 1 easy and fun plus i times 9 root 2 of! Share posts by email math, please use our Google custom search here numbers... 4I ) is a matter of evaluating what is given in rectangular of! − yj $ \begingroup \$ i have attempted this complex number below an. Was covered in topic 43 Step 1 sketch a graph of the following in rectangular form, multiplying. All three four digit numbers formed with non zero digits standard form, the multiplying dividing. See section 2.4 of the text for an introduction to complex numbers drawing. Calculator will simplify any complex expression, with steps shown your Google account some kind of mathematical! D ) Write a multiplying complex numbers in rectangular form number a + jb ; where a complex number from form! Answer is ( 4 ) with a=7, and Last terms together = √-1 work with formulas developed by mathematician... Jb ; where a is the distance from the origin to the point: See.. Sorry, your blog can not share posts by email the standard form, on the other,... On complex numbers ; Graphical explanation of multiplying and adding numbers + . Graphical explanation of multiplying and dividing complex numbers plus i times 9 root 2 over 2 again 18! Imaginary components 2.4 of the number i is defined as i = √-1 multiplying complex numbers in rectangular form 1 and 2. Does basic arithmetic on complex numbers θ 1 and z 2 = r 2 cis 2! In: you are commenting using your Google account symbol ' + ' is treated as vector.... There are two basic forms of complex numbers when they 're in form... Polar coordinates when polar form. sorry, your blog can not share by... Without drawing vectors, we will learn how to multiply complex numbers is in! Sorry, your blog can not share posts by email Euler formula Euler! Numbers that have the form are plotted in the rectangular form. ( a +,. Do you Write a complex number by a real number just as we would with binomial! Determine which of the following development uses trig.formulae you will meet in 36... Followed by 1 unit up on the unit circle need some kind standard. Basic forms of complex numbers and evaluates expressions in the form a complex number we use formulas. The cartesian form of a complex number from polar form. the FOIL method complex Hub aims make... Section `, how to Write the given complex number is written as a+bi where a is called real. Is also called the real axis and the move across and then up convert from polar.! Was covered in topic 43 - 4 learning about complex numbers indices involved., and subtraction of complex numbers without drawing vectors, can also expressed... | 2021-04-15T01:19:12 | {
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https://math.stackexchange.com/questions/2179751/bounded-linear-function-on-ca-b-given-by-fx-xt-0/2179769 | # Bounded linear function on $C[a,b]$ given by $f(x)=x(t_{0})$
I have the following question from a textbook.
Let $f:C[a,\,b]\rightarrow\mathbb{R}$ be defined by $f(x)=x(t_{0})$ for a fixed $t_{0}\in[a,\,b]$. Show that $f$ is a bounded linear function and $\|f\|=1$.
In order to show this I thought I should show that $f$ is continuous and use the following theorem
Theorem: $f$ is continuous if and only if it is bounded. How ever am stuck on how to do this. Is there a way to do this with this technique?
• What norm are you using for f? And by $f(x) = x(t_0)$, is this just multiplication by $t_0$? – Matt Mar 9 '17 at 21:11
• @Matt since the domain of $f$ is the set of continuous functions, $x$ is a function and $x(t_0)$ is the value of $x$ at the point $t_0$, i.e. $f$ could be called evaluation functional. – Roland Mar 9 '17 at 21:35
• @Roland thanks for clarifying! – Matt Mar 9 '17 at 22:37
• It would be good if you could cite the reference of the book! – BAYMAX Sep 14 '17 at 2:32
Lets look at $\|f(x)\|$ and try to express it (or bound) by using $\|x\| = \max_{t\in[a,b]}x(t)$.
$$\|f(x)\| = |x(t_0)| \le \|x\| \Rightarrow \frac{\|f(x)\|}{\|x\|} \le 1$$
But there always exists a constant continous function $x_c(t)$ for which we have $\|x_c\| = x_c(t_0)$ and
$$\|f(x_c)\| = |x_c(t_0)| = \|x_c\| \Rightarrow \frac{\|f(x_c)\|}{\|x_c\|} = 1$$
So, we have $$\|f\| = \sup_{x\in C[a,b]}\frac{\|f(x)\|}{\|x\|} = 1$$
• Won't we need to show that $$|f(x)|\leq M ||x||$$ for some $M\in\mathbb{R}$? so as to show boundedness? – Mafeni Alpha Mar 12 '17 at 13:10
• @MafeniAlpha, the first inequality states exactly this – Andrei Kulunchakov Mar 12 '17 at 13:13
• I have edited the post - the part saying "$\|f(x)\| = x(t_0)$" was missing absolute value - I assume it was just a typo. (Of course, if you prefer to have $\|x(t_0)\|$ rather than $|x(t_0)|$, as to have the notation consistent with what is mostly used for normed spaces, certainly go ahead and edit the post to the form you are satisfied with. – Martin Sleziak Sep 15 '17 at 2:24 | 2020-05-28T16:03:11 | {
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http://math.stackexchange.com/questions/634800/is-2-mathbbz-4-mathbbz-not-a-field/634805 | # Is $2\mathbb{Z}/4\mathbb{Z}$ NOT a field?!
According to wikipedia:
If $R$ is a unital commutative ring with an ideal $m$, then $k = R/m$ is a field if and only if $m$ is a maximal ideal. In that case, $R/m$ is known as the residue field. This fact can fail in non-unital rings. For example, $4\mathbb{Z}$ is a maximal ideal in $2\mathbb{Z}$ , but $2\mathbb{Z}/4\mathbb{Z}$ is not a field.
Is $2\mathbb{Z}/4\mathbb{Z}$ NOT a field?! Isn't it isomorphic to $\mathbb{Z_2}$?
I mean every element in $2\mathbb{Z}$ is of the form $2k$ for some $k \in \mathbb{Z}$. Now suppose that I define $\psi: \mathbb{2Z} \to \mathbb{Z}_2$ where $2k \mapsto \bar{0}$ if $k$ is even and $2k \mapsto \bar{1}$ if $k$ is odd.
This map is well-defined because we have only one representation as $2k$ for each element of $2\mathbb{Z}$. This map is surjective because at least $0$ is mapped to $\bar{0}$ and $2$ is mapped to $\bar{1}$.
Is $\psi$ a homomorphism? Well, to me it sounds like it must be. Because the factor of $2$ is playing no role in here and only the parity of $k$ matters for the mapping so I think it is obvious that it is a homomorphism. Even though it can be checked directly by checking all 4 cases separately as well.
So, what is the kernel? Isn't the kernel equal to $4\mathbb{Z}$? Doesn't that mean $2\mathbb{Z}/4\mathbb{Z}$ is isomorphic to $\mathbb{Z}_2$?
It makes sense that these two are isomorphic to me, where am I wrong?
Also, is the ideal correspondence theorem wrong for non-unital rings???
-
Your map is not a ring homomorphism. If I take $a = 2$, $b = 6$, I have $\psi(2) = \psi(6) = 1$, but $\psi(12) = 0$. Likewise, $2\mathbf{Z}/4\mathbf{Z}$ is not an integral ring.
-
If you allow non-unital rings, you have a nice counterexample here telling you that not every ring is isomorphic to $\mathbf{Z}/2\mathbf{Z}$. :) – fkraiem Jan 11 at 16:30
No. As an abelian group it is indeed isomorphic, but as a ring it needn't be. $2\mathbb{Z}/4\mathbb{Z}$ has only 2 elements, 0 and $\bar{2}$. And $\bar{2}^2=0$, while in $\mathbb{Z}_2$ it isn't. – Giulio Belletti Jan 11 at 16:30
Yeah, I was unaware of the fact that there is more than one ring structure with two elements. Apparently, excluding $\mathbb{Z}_2$, there is at least another ring $R=\{0,1\}$ with the ring multiplication table such that $0.0=0$, $1.0=0$, $0.1=0$, $1.1=0$. I just checked it satisfies the ring axioms. I guess $2\mathbb{Z}/4\mathbb{Z}$ as a ring must be isomorphic to the non-unital ring with two elements that I just mentioned. – some1.new4u Jan 11 at 16:41
In general, given an abelian group, the multiplication defined by $ab = 0$ for all $a,b$ gives a non-unital ring, this is what you have here. – fkraiem Jan 11 at 16:44
@fkraiem: Thanks. Since earlier I had dealt with only commutative rings and it has become a tradition to consider only unital rings in commutative algebra courses I had never seen a non-unital ring before, so this was a delightful surprise for me today. Thanks for your help. – some1.new4u Jan 11 at 16:59
Hint $\ 2\Bbb Z/4\Bbb Z\,$ is not a field since $\,2^2 = 0.\,$ Your map $\psi$ is not a ring hom since they must map the nilpotent $\,2\,$ to a nilpotent ($= 0$ in a field), but you have $\,\psi(2)= 1\,$ by definition. Said in further detail $\ 0 = \psi(0) = \psi(2^2)= \psi(2)^2,\,$ so $\,\psi(2) = 0,\,$ being in a field (or domain).
Remark $\$ Properties like nilpotent, idempotent, invertible etc. that are purely "ring-theoretic" are always preserved by ring homomorphisms. As above, generally this implies constraints on where these elements may be mapped when defining ring homs.
- | 2014-12-21T09:39:41 | {
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https://math.stackexchange.com/questions/3358361/construct-a-function-with-each-derivative-being-non-differentiable-at-a-distinct | # Construct a function with each derivative being non-differentiable at a distinct point
Construct a function $$f_1:\mathbb{R}\to\mathbb{R}$$ with the following properties or show that no such function exists:
$$1.$$ $$f_1$$ is differentiable everywhere except one point $$x_1.$$
$$2.$$ Define $$f_2 : \mathbb{R}\setminus\{x_1\} \to \mathbb{R}$$ as $$f_2(x) :=$$ derivative of $$f_1$$ at $$x.$$ This $$f_2$$ must be differentiable everywhere in its domain except one point $$x_2.$$
$$3.$$ Define $$f_3 : \mathbb{R}\setminus\{x_1,\;x_2\} \to \mathbb{R}$$ as $$f_3(x) :=$$ derivative of $$f_2$$ at $$x.$$ This $$f_3$$ must be differentiable everywhere in its domain except one point $$x_3.$$
$$\vdots$$
$$n.$$ Define $$f_n : \mathbb{R}\setminus\{x_1, \cdots, x_{n-1}\} \to \mathbb{R}$$ as $$f_n(x) :=$$ derivative of $$f_{n-1}$$ at $$x.$$ This $$f_n$$ must be differentiable everywhere in its domain except one point $$x_n.$$
$$\vdots$$
(Note that we do not stop at any $$n.$$)
I found this question in a collection of extra questions for my Calculus course.
I started out by trying something along the lines of $$f(x) = \lim_{n\to\infty}\sum_{i=1}^n (x-i)^i|x-i|$$, but the function itself isn't defined anywhere and I couldn't figure out how to fix it with minimum effort.
So next, I tried out something that might actually be defined somewhere such as $$f(x) = \lim_{n\to\infty}\sum_{i=1}^n \frac{\left(\frac{2}{\pi}\arctan(x-i)\right)^i\left(\frac{2}{\pi}\arctan|x-i|\right)}{(i+1)!}$$
which is defined for $$x \in \mathbb{R}$$, but I wasn't able to prove continuity or differentiability. Intuitively, I feel that since it's a sum of continuous functions, it should be continuous, but I'm not sure whether this intuition is correct because it's an infinite sum.
I asked the person whose website I found the question on (another student), and he said that he wasn't sure whether such a function was even possible.
Any help would be appreciated!
• You were on the right track with your first attempt (but replace $x-i$ by $x-x_i$, where $(x_i)$ is the given sequence of points). It is true that the series needs not converge at any point, but that can be remedied easily by multiplying every summand by a sequence of numbers that decreases sufficiently fast, such as $i!$... – Giuseppe Negro Sep 16 at 12:35
A different solution is the following 'lazy exponential' - there are easier solutions (maybe look up bump functions), but I like delay ODEs. set \begin{align}x\in(-\infty,0]&\implies f(x):=1,\\ x\in(0,1] &\implies f(x) := 1+x, \\ x\in (1,2] &\implies f(x) := 1+x + \frac{(x-1)^2}{2!},\\ x\in(2,3] & \implies f(x) := 1+x + \frac{(x-1)^2}{2!} + \frac{(x-2)^3}{3!}, \end{align} and in general $$x\in(n,n+1]\implies f(x) := \sum_{k=0}^{n+1} \frac{(x-k+1)^k}{k!}.$$
If you differentiate, you find for $$x\in (n,n+1)$$, where $$n>1$$:
$$f'(x) = \sum_{k=1}^{n+1} \frac{(x-k+1)^{k-1}}{(k-1)!} =\sum_{j=0}^{n} \frac{(x-j)^{j}}{j!}= f(x-1)$$ so to the right of 1, it solves a delay ODE with initial data prescribed on $$x\in(0,1]$$ above. $$f'$$ is clearly discontinuous at $$0$$, but $$\left.\frac{d}{dx}\frac{(x-1)^2}{2!}\right|_{x=1} = 0$$ so the derivative is continuous at $$x=1$$. In general, for any integer $$n\ge 2$$, near $$x=n-1$$, all the terms $$\frac{(x-h+1)^h}{h!}$$ for $$h are smooth, and the newly added term $$T_n$$, $$T_n(x) := \begin{cases} \frac{(x-n+1)^n}{n!} & x>n-1,\\ 0 & x\le n-1\end{cases}$$ is $$C^1$$. Conclusion - $$f \in C^0(\mathbb R)\cap C^1(\mathbb R\setminus \{0\}).$$
To finish, we use the delay ODE, which says that differentiating is the same as translating the function to the right by one. Thus for $$x\in \mathbb (0,\infty)\setminus \mathbb N$$, $$i\in\mathbb N$$, $$f^{(i+1)}(x+i) = f'(x).$$ So the discontinuity of $$f^{(i+1)}$$ at $$x=i-1$$, and the continuity at integers $$x=\tilde i > i-1$$ follows directly from the dis/continuity of $$f'$$ at $$0,1,2,\dots$$. We conclude $$f \in C^0(\mathbb R)\cap \left(\bigcap_{k=1}^\infty C^k(\mathbb R\setminus{\{0,1,\dots,k-1\}})\right).$$ | 2019-12-13T16:24:27 | {
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https://mathhelpboards.com/threads/need-help-with-parabolas.5786/ | # Need help with parabolas
##### New member
I've encountered a question that I need help understanding the answer to.
The question is:
What is the equation of the axis of symmetry of the parabola given by the equation y=6(x+1)(x-5)?
Now, I know this quadratic equation is in intercept form, and I know that the formula for finding the axis of symmetry for this is (p+q)/2
Which means that it's (1-5)/2
then -4/2
Which equals -2. But it says that the correct answer is 2, not -2.
#### Bacterius
##### Well-known member
MHB Math Helper
[JUSTIFY]That's because if $(x - a)$ is a factor of your equation, then the intercept is at $x = a$, not $x = -a$. So in fact your parabola has intercepts $x = -1$ and $x = 5$, and so you would get $\frac{-1 + 5}{2} = \frac{4}{2} = 2$ as expected. The formula itself is correct.
As an example, try plotting $y = x - 5$, and you'll see it intercepts the x-axis at $x = 5$, and not $x = -5$. In the same way, try plotting your parabola and see where it intercepts the x-axis.[/JUSTIFY] | 2021-01-24T11:30:43 | {
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http://math.stackexchange.com/questions/169407/find-identity-element-invertible-and-inverses-in-t-mathbb-z-times-mathbb-q | # Find identity element, invertible and inverses in $T=\mathbb Z \times \mathbb Q$
Let the following operation be defined on $T=\mathbb Z \times \mathbb Q$:
\begin{aligned}(a,b)\centerdot (c,d) = (-ac, b+d+2) \end{aligned}
in the commutative semigroup $(T, \centerdot)$, find the identity element, invertible elements and inverses.
I need to find the identity element, so I need an $(\alpha, \beta)$ so that
\begin{aligned}(a,b)\centerdot (\alpha,\beta) = (a,b) \end{aligned}
as $(T, \centerdot)$ is a commutative semigroup, I assume that $(\alpha,\beta)\centerdot (a,b) = (a,b)$ holds.
\begin{aligned}-a\alpha = a \Leftrightarrow \alpha = -1 \end{aligned} \begin{aligned}b+\beta+2 = b \Leftrightarrow \beta = -2 \end{aligned}
so the identity element relative to this semigroup is $(-1,-2)$.
In order to find the invertible elements and inverses:
\begin{aligned}(a,b)\centerdot (a',b') = (-1,-2) \end{aligned}
so
\begin{aligned}-aa' = -1 \Leftrightarrow a'=-\frac{1}{a} \in \mathbb Z \Rightarrow a = 1 \end{aligned}
\begin{aligned}b+b'+2=-2 \Leftrightarrow b'=-(b+4) \end{aligned}
in conclusion only the elements $(1,b)$ are actually invertible and their inverse is $(1, -(b+4))$.
Does this whole thing hold? Am I wrong in any part of it?
-
In the last bit, why can't $a$ be $-1$? – Gerry Myerson Jul 11 '12 at 9:37
@GerryMyerson because I need $(-aa') = -1$ if $a = -1$ then $(-aa') = 1$. – haunted85 Jul 11 '12 at 9:42
As @GerryMyerson pointed out $a$ cab be $-1\implies a'=-1\implies$ inverse of $(-1,b)$ also exists and is equal to $(-1,-(b+4)).$ – Aang Jul 11 '12 at 9:46
After equation $-aa'=-1$ $a'=\frac{-1}{a}$ is wrong, it should be $\frac{1}{a}$. – Aang Jul 11 '12 at 9:49
@avatar thank you I've really overlooked that minus sign!! what a shame! – haunted85 Jul 11 '12 at 9:59
In all of these problems, you are taking a structure with operations that you "already know", and are trying to test to see whether a new operation, defined in terms of the ones you know, satisfies certain properties.
It's usually best to denote the "new" operation with a symbol that is unlikely to generate confusion. You've been doing this elsewhere, but here you are using $\cdot$, which can easily be confused with regular integer multiplication. So I'm going to replace it with $\odot$, if you don't mind.
Let's approach it systematically, and essentially "following your nose." I'm going to check a few things that you seem to take for granted, just to get familiarity with the operation.
Now, we have the set $T=\mathbb{Z}\times\mathbb{Q}$, and the operation $$(a,b)\odot (c,d) = (-ac, b+d+2).$$ First, the operation is commutative, since integer multiplication and addition of rationals is commutative. It is also associative: \begin{align*} \Bigl( (a,b)\odot (c,d)\Bigr)\odot (x,y) &= (-ac, b+d+2)\odot (x,y)\\ &=(-(-ac)x,(b+d+2)+y+2)\\ &= (acx, b+d+y+4)\\ (a,b)\odot \Bigl((c,d)\odot (x,y)\Bigr)&= (a,b)\odot (-cx, d+y+2)\\ &= (-a(-cx), b+(d+y+2)+2)\\ &= (acx, b+d+y+4). \end{align*} So the operation is commutative and associative, and we do indeed have a semigroup.
Now we are trying to see whether this is a monoid (semigroup with identity). So we are trying to see whether there exist $\alpha\in\mathbb{Z}$ and $\beta\in\mathbb{Q}$ such that, for all $a\in\mathbb{Z}$ and $b\in\mathbb{Q}$, we have $$(\alpha,\beta)\odot(a,b) = (a,b)$$ (exactly as you did). This amounts to solving the equations you solved, which you did correctly, and so we discovered that $(T,\odot)$ is indeed a monoid, and that the identity element is $e_T = (-1,-2)$. Great!
To figure out what elements have inverses, you do set $(a,b)\odot (x,y) =e_T = (-1,-2)$, and figure out if you can determine necessary and sufficient conditions on $a$ and $b$ for $x$ and $y$ to exist. As you note, we have: $$(-1,-2) = (a,b)\odot(x,y) = (-ax, b+y+2).$$ This requires $-ax = 1$, which means that $a$ has to be an integer that has a multiplicative inverse in the integers (namely, $-x$); there are only two possibilities: $a=1$ or $a=-1$.
On the other hand, for any rational number $b$ we can always find a rational number $y$ such that $b+y+2=-2$, namely, as you note, $y = -(b+4)$. So it would seem that provided that $a=\pm 1$, then $(a,b)$ has an $\odot$-inverse in $T$. And, indeed, if $a=\pm 1$, then we have $$(\pm 1,b)\odot (\pm1,-b-4) = (-(\pm1)(\pm 1), b-b-4+2) = (-1-2).$$ So a necessary and sufficient condition for $(a,b)\in T$ to be $\odot$-invertible is that $a=\pm 1$.
-
thank you so much for your patient, clear and most of all exhaustive answer, appreciated. – haunted85 Jul 11 '12 at 19:23 | 2016-06-25T14:17:24 | {
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https://www.physicsforums.com/threads/derivative-of-x-to-the-x.155954/ | Derivative of x to the x
1. Feb 12, 2007
lapenna
1. The problem statement, all variables and given/known data
I am required to find the derivative of x^x using the differance quotient
2. Relevant equations
((x+h)^(x+h)-(x^x))/h
3. The attempt at a solution
TI-89 gives (lnx+1)x^x I can't figure out how to get there
2. Feb 12, 2007
theperthvan
Without the derivative definition, ln both sides, so:
$$y=x^x$$
$$ln(y)=xln(x)$$
$$\frac{1}{y}.\frac{dy}{dx} = 1.ln(x) + x.\frac{1}{x}$$
$$\frac{dy}{dx} = y[ln(x) + 1]$$
$$\frac{dy}{dx} = x^x(ln(x) + 1)$$
3. Feb 13, 2007
HallsofIvy
Staff Emeritus
I can't imagine there being any reasonable way to use the "difference quotient" to determine the derivative of xx. You must have a really evil teacher!
4. Feb 14, 2007
dextercioby
Either this is unreasonable, or your imagination really needs to see more:
$$x^x =e^{x\ln x}$$
$$\lim_{h\rightarrow 0} \frac{(x+h)^{(x+h)}-x^x}{h} =\lim_{h\rightarrow 0} \frac{e^{(x+h)\ln(x+h)}-e^{x\ln x}}{h}=\lim_{h\rightarrow 0}\frac{e^{(x+h)\ln\left(x\left(1+\frac{h}{x}\right)\right)}-e^{x\ln x}}{h}$$
$$=\lim_{h\rightarrow 0} \frac{e^{x\ln x}e^{h\ln x}e^{x\ln\left(1+\frac{h}{x}\right)}e^{h\ln\left(1+\frac{h}{x}\right)}-e^{x\ln x}}{h}$$
$$=x^{x}\lim_{h\rightarrow 0}\frac{e^{h\ln x}e^{x\ln\left(1+\frac{h}{x}\right)}e^{h\ln\left(1+\frac{h}{x}\right)}-1}{h}\cdot \frac{h\ln x+x\ln\left(1+\frac{h}{x}\right)+h\ln\left(1+\frac{h}{x}\right)}{h\ln x+x\ln\left(1+\frac{h}{x}\right)+h\ln\left(1+\frac{h}{x}\right)}$$
$$=x^{x}\left[\ln x+x\lim_{h\rightarrow 0}\frac{\ln\left(1+\frac{h}{x}\right)}{h}+\lim_{h\rightarrow 0}\ln\left(1+\frac{h}{x}\right)\right]$$
$$=x^x \left\{\ln x+x\ln \lim_{h\rightarrow 0}\left[\left(1+\frac{h}{x}\right)^{\frac{x}{h}}\right]^{\frac{1}{x}}\right\}$$
$$=x^x \left(\ln x +1\right)$$
Last edited: Feb 14, 2007
5. Feb 15, 2007
theperthvan
wow .
6. Feb 15, 2007
VietDao29
Whoa.
Or somewhat easier. That's a little bit messy methinks:
$$(x ^ x)' = (e ^ {x \ln (x)})' = \lim_{h \rightarrow 0} \frac{e ^ {(x + h) \ln (x + h)} - e ^ {x \ln (x)}}{h}$$
$$= \lim_{h \rightarrow 0} \left( \frac{e ^ {(x + h) \ln (x + h)} - e ^ {x \ln (x)}}{(x + h) \ln (x + h) - x \ln x} \times \frac{(x + h) \ln (x + h) - x \ln x}{h} \right)$$
$$= e ^ {x \ln (x)} \left( \lim_{h \rightarrow 0} \frac{x \ln (1 + \frac{h}{x}) + h \ln (x + h)}{h} \right)$$ [due to: (eu)'u = eu]
$$= x ^ x \left\{ \lim_{h \rightarrow 0} \left[ x \ln \left( \left( 1 + \frac{h}{x} \right) ^ {\frac{x}{h}} \right) ^ {\frac{1}{x}} \right] + \ln (x) \right\}$$
$$= x ^ x (1 + \ln (x))$$
Last edited: Feb 15, 2007
7. Aug 21, 2009
Raiey
Since the original equation is x^x and never a division quotient, one can not apply that rule. Here is the alternative proof.
Attached Files:
• derivativesxtopowerofx.txt
File size:
3.4 KB
Views:
45
8. Aug 22, 2009
g_edgar
"Difference quotient" means this:
$$\frac{f(x)-f(a)}{x-a}$$
so the assignment was to use the definition of derivative that involves this.
9. Aug 23, 2009
Raiey
An answer is attached in PDF
File size:
37.7 KB
Views:
78 | 2017-10-17T02:12:13 | {
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https://math.libretexts.org/Bookshelves/Precalculus/Book%3A_Precalculus_(Stitz-Zeager)/10%3A_Foundations_of_Trigonometry/10.4%3A_Trigonometric_Identities |
# 10.4: Trigonometric Identities
In Section \ref{CircularFunctions}, we saw the utility of the Pythagorean Identities in Theorem \ref{pythids} along with the Quotient and Reciprocal Identities in Theorem \ref{recipquotid}. Not only did these identities help us compute the values of the circular functions for angles, they were also useful in simplifying expressions involving the circular functions. In this section, we introduce several collections of identities which have uses in this course and beyond. Our first set of identities is the Even / Odd' identities.\footnote{As mentioned at the end of Section \ref{TheUnitCircle}, properties of the circular functions when thought of as functions of angles in radian measure hold equally well if we view these functions as functions of real numbers. Not surprisingly, the Even / Odd properties of the circular functions are so named because they identify cosine and secant as even functions, while the remaining four circular functions are odd. (See Section \ref{GraphsofFunctions}.)}
Note: Even / Odd Identities
For all applicable angles $$\theta$$:
• $$\cos(-\theta) = \cos(\theta)$$
• $$\sec(-\theta) = \sec(\theta)$$
• $$\sin(-\theta) = -\sin(\theta)$$
• $$\csc(-\theta) = -\csc(\theta)$$
• $$\tan(-\theta) = -\tan(\theta)$$
• $$\cot(-\theta) = -\cot(\theta)$$
In light of the Quotient and Reciprocal Identities, Theorem \ref{recipquotid}, it suffices to show $$\cos(-\theta) = \cos(\theta)$$ and $$\sin(-\theta) = -\sin(\theta)$$. The remaining four circular functions can be expressed in terms of $$\cos(\theta)$$ and $$\sin(\theta)$$ so the proofs of their Even / Odd Identities are left as exercises. Consider an angle $$\theta$$ plotted in standard position. Let $$\theta_ { o}$$ be the angle coterminal with $$\theta$$ with $$0 \leq \theta_ { o} < 2\pi$$. (We can construct the angle $$\theta_ { o}$$ by rotating counter-clockwise from the positive $$x$$-axis to the terminal side of $$\theta$$ as pictured below.) Since $$\theta$$ and $$\theta_ { o}$$ are coterminal, $$\cos(\theta) = \cos(\theta_ { o})$$ and $$\sin(\theta) = \sin(\theta_ { o})$$.
We now consider the angles $$-\theta$$ and $$-\theta_ { o}$$. Since $$\theta$$ is coterminal with $$\theta_ { o}$$, there is some integer $$k$$ so that $$\theta = \theta_ { o} + 2\pi \cdot k$$. Therefore, $$-\theta = -\theta_ { o} - 2\pi \cdot k = -\theta_ { o} + 2\pi \cdot(-k)$$. Since $$k$$ is an integer, so is $$(-k)$$, which means $$-\theta$$ is coterminal with $$-\theta_ { o}$$. Hence, $$\cos(-\theta) = \cos(-\theta_ { o})$$ and $$\sin(-\theta) = \sin(-\theta_ { o})$$. Let $$P$$ and $$Q$$ denote the points on the terminal sides of $$\theta_ { o}$$ and $$-\theta_ { o}$$, respectively, which lie on the Unit Circle. By definition, the coordinates of $$P$$ are $$(\cos(\theta_ { o}),\sin(\theta_ { o}))$$ and the coordinates of $$Q$$ are $$(\cos(-\theta_ { o}),\sin(-\theta_ { o}))$$. Since $$\theta_ { o}$$ and $$-\theta_ { o}$$ sweep out congruent central sectors of the Unit Circle, it follows that the points $$P$$ and $$Q$$ are symmetric about the $$x$$-axis. Thus, $$\cos(-\theta_ { o}) = \cos(\theta_ { o})$$ and $$\sin(-\theta_ { o}) = -\sin(\theta_ { o})$$. Since the cosines and sines of $$\theta_ { o}$$ and $$-\theta_ { o}$$ are the same as those for $$\theta$$ and $$-\theta$$, respectively, we get $$\cos(-\theta) = \cos(\theta)$$ and $$\sin(-\theta) = -\sin(\theta)$$, as required. The Even / Odd Identities are readily demonstrated using any of the common angles' noted in Section \ref{TheUnitCircle}. Their true utility, however, lies not in computation, but in simplifying expressions involving the circular functions. In fact, our next batch of identities makes heavy use of the Even / Odd Identities.
Note: Sum and Difference Identities for Cosine
For all angles $$\alpha$$ and $$\beta$$:
• $$\cos(\alpha + \beta) = \cos(\alpha) \cos(\beta) - \sin(\alpha) \sin(\beta)$$
• $$\cos(\alpha - \beta) = \cos(\alpha) \cos(\beta) + \sin(\alpha) \sin(\beta)$$
We first prove the result for differences. As in the proof of the Even / Odd Identities, we can reduce the proof for general angles $$\alpha$$ and $$\beta$$ to angles $$\alpha_ { o}$$ and $$\beta_ { o}$$, coterminal with $$\alpha$$ and $$\beta$$, respectively, each of which measure between $$0$$ and $$2\pi$$ radians. Since $$\alpha$$ and $$\alpha_ { o}$$ are coterminal, as are $$\beta$$ and $$\beta_ { o}$$, it follows that $$\alpha - \beta$$ is coterminal with $$\alpha_ { o} - \beta_ { o}$$. Consider the case below where $$\alpha_ { o} \geq \beta_ { o}$$.
Since the angles $$POQ$$ and $$AOB$$ are congruent, the distance between $$P$$ and $$Q$$ is equal to the distance between $$A$$ and $$B$$.\footnote{In the picture we've drawn, the \underline{tri}angles $$POQ$$ and $$AOB$$ are congruent, which is even better. However, $$\alpha_ { o} - \beta_ { o}$$ could be $$0$$ or it could be $$\pi$$, neither of which makes a triangle. It could also be larger than $$\pi$$, which makes a triangle, just not the one we've drawn. You should think about those three cases.} The distance formula, Equation \ref{distanceformula}, yields
$\begin{array}{rcl} \sqrt{(\cos(\alpha_ { o}) - \cos(\beta_ { o}))^2 + (\sin(\alpha_ { o}) - \sin(\beta_ { o}))^2 } & = & \sqrt{(\cos(\alpha_ { o} - \beta_ { o}) - 1)^2 + (\sin(\alpha_ { o} - \beta_ { o}) - 0)^2} \\ \end{array}$
Squaring both sides, we expand the left hand side of this equation as
$\begin{array}{rcl} (\cos(\alpha_ { o}) - \cos(\beta_ { o}))^2 + (\sin(\alpha_ { o}) - \sin(\beta_ { o}))^2 & = & \cos^2(\alpha_ { o}) - 2\cos(\alpha_ { o})\cos(\beta_ { o}) + \cos^2(\beta_ { o}) \\ & & + \sin^2(\alpha_ { o}) - 2\sin(\alpha_ { o})\sin(\beta_ { o}) + \sin^2(\beta_ { o}) \\ & = & \cos^2(\alpha_ { o}) + \sin^2(\alpha_ { o}) + \cos^2(\beta_ { o}) + \sin^2(\beta_ { o}) \\ & & - 2\cos(\alpha_ { o})\cos(\beta_ { o}) - 2\sin(\alpha_ { o})\sin(\beta_ { o}) \end{array}$
From the Pythagorean Identities, $$\cos^2(\alpha_ { o}) + \sin^2(\alpha_ { o}) = 1$$ and $$\cos^2(\beta_ { o}) + \sin^2(\beta_ { o}) = 1$$, so
$\begin{array}{rcl} (\cos(\alpha_ { o}) - \cos(\beta_ { o}))^2 + (\sin(\alpha_ { o}) - \sin(\beta_ { o}))^2 & = & 2 - 2\cos(\alpha_ { o})\cos(\beta_ { o}) - 2\sin(\alpha_ { o})\sin(\beta_ { o}) \end{array}$
Turning our attention to the right hand side of our equation, we find
$\begin{array}{rcl} (\cos(\alpha_ { o} - \beta_ { o}) - 1)^2 + (\sin(\alpha_ { o} - \beta_ { o}) - 0)^2 & = & \cos^2(\alpha_ { o} - \beta_ { o}) - 2\cos(\alpha_ { o} - \beta_ { o}) + 1 + \sin^2(\alpha_ { o} - \beta_ { o}) \\ & = & 1 + \cos^2(\alpha_ { o} - \beta_ { o}) + \sin^2(\alpha_ { o} - \beta_ { o}) - 2\cos(\alpha_ { o} - \beta_ { o}) \\ \end{array}$
Once again, we simplify $$\cos^2(\alpha_ { o} - \beta_ { o}) + \sin^2(\alpha_ { o} - \beta_ { o})= 1$$, so that
$\begin{array}{rcl} (\cos(\alpha_ { o} - \beta_ { o}) - 1)^2 + (\sin(\alpha_ { o} - \beta_ { o}) - 0)^2 & = & 2 - 2\cos(\alpha_ { o} - \beta_ { o}) \\ \end{array}$
Putting it all together, we get $$2 - 2\cos(\alpha_ { o})\cos(\beta_ { o}) - 2\sin(\alpha_ { o})\sin(\beta_ { o}) = 2 - 2\cos(\alpha_ { o} - \beta_ { o})$$, which simplifies to: $$\cos(\alpha_ { o} - \beta_ { o}) = \cos(\alpha_ { o})\cos(\beta_ { o}) + \sin(\alpha_ { o})\sin(\beta_ { o})$$. Since $$\alpha$$ and $$\alpha_ { o}$$, $$\beta$$ and $$\beta_ { o}$$ and $$\alpha - \beta$$ and $$\alpha_ { o}- \beta_ { o}$$ are all coterminal pairs of angles, we have $$\cos(\alpha - \beta) = \cos(\alpha) \cos(\beta) + \sin(\alpha) \sin(\beta)$$. For the case where $$\alpha_ { o} \leq \beta_ { o}$$, we can apply the above argument to the angle $$\beta_ { o} - \alpha_ { o}$$ to obtain the identity $$\cos(\beta_ { o} - \alpha_ { o}) = \cos(\beta_ { o})\cos(\alpha_ { o}) + \sin(\beta_ { o})\sin(\alpha_ { o})$$. Applying the Even Identity of cosine, we get $$\cos(\beta_ { o} - \alpha_ { o}) = \cos( - (\alpha_ { o} - \beta_ { o})) = \cos(\alpha_ { o} - \beta_ { o})$$, and we get the identity in this case, too.
To get the sum identity for cosine, we use the difference formula along with the Even/Odd Identities
$\cos(\alpha + \beta) = \cos(\alpha - (-\beta)) = \cos(\alpha) \cos(-\beta) + \sin(\alpha) \sin(-\beta) = \cos(\alpha) \cos(\beta) - \sin(\alpha) \sin(\beta)$
We put these newfound identities to good use in the following example.
Example $$\PageIndex{1}$$: Cosine Sum and Difference
1. Find the exact value of $$\cos\left(15^{\circ}\right)$$.
2. Verify the identity: $$\cos\left(\frac{\pi}{2} - \theta\right) = \sin(\theta)$$.
Solution
1. In order to use Theorem \ref{cosinesumdifference} to find $$\cos\left(15^{\circ}\right)$$, we need to write $$15^{\circ}$$ as a sum or difference of angles whose cosines and sines we know. One way to do so is to write $$15^{\circ} = 45^{\circ} - 30^{\circ}$$.
$\begin{array}{rcl} \cos\left(15^{\circ}\right) & = & \cos\left(45^{\circ} - 30^{\circ} \right) \\ & = & \cos\left(45^{\circ}\right)\cos\left(30^{\circ} \right) + \sin\left(45^{\circ}\right)\sin\left(30^{\circ} \right) \\ & = & \left( \dfrac{\sqrt{2}}{2} \right)\left( \dfrac{\sqrt{3}}{2} \right) + \left( \dfrac{\sqrt{2}}{2} \right)\left( \dfrac{1}{2} \right)\\ & = & \dfrac{\sqrt{6}+ \sqrt{2}}{4} \\ \end{array}$
1. In a straightforward application of Theorem \ref{cosinesumdifference}, we find
$\begin{array}{rcl} \cos\left(\dfrac{\pi}{2} - \theta\right) & = & \cos\left(\dfrac{\pi}{2}\right)\cos\left(\theta\right) + \sin\left(\dfrac{\pi}{2}\right)\sin\left(\theta \right) \\ & = & \left( 0 \right)\left( \cos(\theta) \right) + \left( 1 \right)\left( \sin(\theta) \right) \\ & = & \sin(\theta) \\ \end{array}$
The identity verified in Example $$\PageIndex{1}$$, namely, $$\cos\left(\frac{\pi}{2} - \theta\right) = \sin(\theta)$$, is the first of the celebrated cofunction' identities. These identities were first hinted at in Exercise \ref{cofunctionforeshadowing} in Section \ref{TheUnitCircle}. From $$\sin(\theta) = \cos\left(\frac{\pi}{2} - \theta\right)$$, we get:
$\sin\left(\dfrac{\pi}{2} - \theta\right) = \cos\left(\dfrac{\pi}{2} -\left[\dfrac{\pi}{2} - \theta\right]\right) = \cos(\theta),$
which says, in words, that the co'sine of an angle is the sine of its co'mplement. Now that these identities have been established for cosine and sine, the remaining circular functions follow suit. The remaining proofs are left as exercises.
Note: Cofunction Identities
For all applicable angles $$\theta$$:
• $$\cos\left(\dfrac{\pi}{2} - \theta \right) = \sin(\theta)$$
• $$\sin\left(\dfrac{\pi}{2} - \theta \right) = \cos(\theta)$$
• $$\sec\left(\dfrac{\pi}{2} - \theta \right) = \csc(\theta)$$
• $$\csc\left(\dfrac{\pi}{2} - \theta \right) = \sec(\theta)$$
• $$\tan\left(\dfrac{\pi}{2} - \theta \right) = \cot(\theta)$$
• $$\cot\left(\dfrac{\pi}{2} - \theta \right) = \tan(\theta)$$
With the Cofunction Identities in place, we are now in the position to derive the sum and difference formulas for sine. To derive the sum formula for sine, we convert to cosines using a cofunction identity, then expand using the difference formula for cosine
$\begin{array}{rcl} \sin(\alpha + \beta) & = & \cos\left( \dfrac{\pi}{2} - (\alpha + \beta) \right) \\ & = & \cos\left( \left[\dfrac{\pi}{2} - \alpha \right] - \beta \right) \\ & = & \cos\left(\dfrac{\pi}{2} - \alpha \right) \cos(\beta) + \sin\left(\dfrac{\pi}{2} - \alpha \right)\sin(\beta) \\ & = & \sin(\alpha) \cos(\beta) + \cos(\alpha) \sin(\beta) \\ \end{array}$
We can derive the difference formula for sine by rewriting $$\sin(\alpha - \beta)$$ as $$\sin(\alpha + (-\beta))$$ and using the sum formula and the Even / Odd Identities. Again, we leave the details to the reader.
Sum and Difference Identities for Sine
For all angles $$\alpha$$ and $$\beta$$, \index{Difference Identity ! for sine} \index{Sum Identity ! for sine}
• $$\sin(\alpha + \beta) = \sin(\alpha) \cos(\beta) + \cos(\alpha) \sin(\beta)$$
• $$\sin(\alpha - \beta) = \sin(\alpha) \cos(\beta) - \cos(\alpha) \sin(\beta)$$
Example $$\PageIndex{1}$$:
1. Find the exact value of $$\sin\left(\frac{19 \pi}{12}\right)$$
2. If $$\alpha$$ is a Quadrant II angle with $$\sin(\alpha) = \frac{5}{13}$$, and $$\beta$$ is a Quadrant III angle with $$\tan(\beta) = 2$$, find $$\sin(\alpha - \beta)$$.
3. Derive a formula for $$\tan(\alpha + \beta)$$ in terms of $$\tan(\alpha)$$ and $$\tan(\beta)$$.
Solution
1. As in Example \ref{cosinesumdiffex}, we need to write the angle $$\frac{19 \pi}{12}$$ as a sum or difference of common angles. The denominator of $$12$$ suggests a combination of angles with denominators $$3$$ and $$4$$. One such combination is $$\; \frac{19 \pi}{12} = \frac{4 \pi}{3} + \frac{\pi}{4}$$. Applying Theorem \ref{sinesumdifference}, we get
$\begin{array}{rcl} \sin\left(\dfrac{19 \pi}{12}\right) & = & \sin\left(\dfrac{4 \pi}{3} + \dfrac{\pi}{4} \right) \\ & = & \sin\left(\dfrac{4 \pi}{3} \right)\cos\left(\dfrac{\pi}{4} \right) + \cos\left(\dfrac{4 \pi}{3} \right)\sin\left(\dfrac{\pi}{4} \right) \\ & = & \left( -\dfrac{\sqrt{3}}{2} \right)\left( \dfrac{\sqrt{2}}{2} \right) + \left( -\dfrac{1}{2} \right)\left( \dfrac{\sqrt{2}}{2} \right) \\ & = & \dfrac{-\sqrt{6}- \sqrt{2}}{4} \\ \end{array}$
1. In order to find $$\sin(\alpha - \beta)$$ using Theorem \ref{sinesumdifference}, we need to find $$\cos(\alpha)$$ and both $$\cos(\beta)$$ and $$\sin(\beta)$$. To find $$\cos(\alpha)$$, we use the Pythagorean Identity $$\cos^2(\alpha) + \sin^2(\alpha) = 1$$. Since $$\sin(\alpha) = \frac{5}{13}$$, we have $$\cos^{2}(\alpha) + \left(\frac{5}{13}\right)^2 = 1$$, or $$\cos(\alpha) = \pm \frac{12}{13}$$. Since $$\alpha$$ is a Quadrant II angle, $$\cos(\alpha) = -\frac{12}{13}$$. We now set about finding $$\cos(\beta)$$ and $$\sin(\beta)$$. We have several ways to proceed, but the Pythagorean Identity $$1 + \tan^{2}(\beta) = \sec^{2}(\beta)$$ is a quick way to get $$\sec(\beta)$$, and hence, $$\cos(\beta)$$. With $$\tan(\beta) = 2$$, we get $$1 + 2^2 = \sec^{2}(\beta)$$ so that $$\sec(\beta) = \pm \sqrt{5}$$. Since $$\beta$$ is a Quadrant III angle, we choose $$\sec(\beta) = -\sqrt{5}$$ so $$\cos(\beta) = \frac{1}{\sec(\beta)} = \frac{1}{-\sqrt{5}} = -\frac{\sqrt{5}}{5}$$. We now need to determine $$\sin(\beta)$$. We could use The Pythagorean Identity $$\cos^{2}(\beta) + \sin^{2}(\beta) = 1$$, but we opt instead to use a quotient identity. From $$\tan(\beta) = \frac{\sin(\beta)}{\cos(\beta)}$$, we have $$\sin(\beta) = \tan(\beta) \cos(\beta)$$ so we get $$\sin(\beta) = (2) \left( -\frac{\sqrt{5}}{5}\right) = - \frac{2 \sqrt{5}}{5}$$. We now have all the pieces needed to find $$\sin(\alpha - \beta)$$:
$\begin{array}{rcl} \sin(\alpha - \beta) & = & \sin(\alpha)\cos(\beta) - \cos(\alpha)\sin(\beta) \\ & = & \left( \dfrac{5}{13} \right)\left( -\dfrac{\sqrt{5}}{5} \right) - \left( -\dfrac{12}{13} \right)\left( - \dfrac{2 \sqrt{5}}{5} \right) \\ & = & -\dfrac{29\sqrt{5}}{65} \\ \end{array}$
We can start expanding $$\tan(\alpha + \beta)$$ using a quotient identity and our sum formulas
$\begin{array}{rcl} \tan(\alpha + \beta) & = & \dfrac{\sin(\alpha + \beta)}{\cos(\alpha + \beta)} \\ & = & \dfrac{\sin(\alpha) \cos(\beta) + \cos(\alpha) \sin(\beta)}{\cos(\alpha) \cos(\beta) - \sin(\alpha) \sin(\beta)} \\ \end{array}$
Since $$\tan(\alpha) = \frac{\sin(\alpha)}{\cos(\alpha)}$$ and $$\tan(\beta) = \frac{\sin(\beta)}{\cos(\beta)}$$, it looks as though if we divide both numerator and denominator by $$\cos(\alpha) \cos(\beta)$$ we will have what we want
$\begin{array}{rcl} \tan(\alpha + \beta) & = & \dfrac{\sin(\alpha) \cos(\beta) + \cos(\alpha) \sin(\beta)}{\cos(\alpha) \cos(\beta) - \sin(\alpha) \sin(\beta)} \cdot\dfrac{\dfrac{1}{\cos(\alpha) \cos(\beta)}}{\dfrac{1}{\cos(\alpha) \cos(\beta)}}\\ & & \\ & = & \dfrac{\dfrac{\sin(\alpha) \cos(\beta)}{\cos(\alpha) \cos(\beta)} + \dfrac{\cos(\alpha) \sin(\beta)}{\cos(\alpha) \cos(\beta)}}{\dfrac{\cos(\alpha) \cos(\beta)}{\cos(\alpha) \cos(\beta)} - \dfrac{\sin(\alpha) \sin(\beta)}{\cos(\alpha) \cos(\beta)}}\\ & & \\ & = & \dfrac{\dfrac{\sin(\alpha) \cancel{\cos(\beta)}}{\cos(\alpha) \cancel{\cos(\beta)}} + \dfrac{\cancel{\cos(\alpha)} \sin(\beta)}{\cancel{\cos(\alpha)} \cos(\beta)}}{\dfrac{\cancel{\cos(\alpha)} \cancel{\cos(\beta)}}{\cancel{\cos(\alpha)} \cancel{\cos(\beta)}} - \dfrac{\sin(\alpha) \sin(\beta)}{\cos(\alpha) \cos(\beta)}}\\ & & \\ & = & \dfrac{\tan(\alpha) + \tan(\beta)}{1 -\tan(\alpha) \tan(\beta)}\\ \end{array}$
Naturally, this formula is limited to those cases where all of the tangents are defined.\qed
The formula developed in Exercise \ref{sinesumanddiffex} for $$\tan(\alpha + \beta)$$ can be used to find a formula for $$\tan(\alpha - \beta)$$ by rewriting the difference as a sum, $$\tan(\alpha + (-\beta))$$, and the reader is encouraged to fill in the details. Below we summarize all of the sum and difference formulas for cosine, sine and tangent.
Note
Sum and Difference Identities:} For all applicable angles $$\alpha$$ and $$\beta$$, \index{Difference Identity ! for tangent} \index{Sum Identity ! for tangent} \index{Difference Identity ! for cosine} \index{Sum Identity ! for cosine} \index{Difference Identity ! for sine} \index{Sum Identity ! for sine}
• $$\cos(\alpha \pm \beta) = \cos(\alpha) \cos(\beta) \mp \sin(\alpha) \sin(\beta)$$
• $$\sin(\alpha \pm \beta) = \sin(\alpha) \cos(\beta) \pm \cos(\alpha) \sin(\beta)$$
• $$\tan(\alpha \pm \beta) = \dfrac{\tan(\alpha) \pm \tan(\beta)}{1 \mp \tan(\alpha) \tan(\beta)}$$
In the statement of Theorem \ref{circularsumdifference}, we have combined the cases for the sum $+$' and difference \)-$' of angles into one formula. The convention here is that if you want the formula for the sum $+$' of two angles, you use the top sign in the formula; for the difference, \)-$', use the bottom sign. For example, $\tan(\alpha - \beta) = \dfrac{\tan(\alpha) - \tan(\beta)}{1 + \tan(\alpha) \tan(\beta)}$
If we specialize the sum formulas in Theorem \ref{circularsumdifference} to the case when $$\alpha = \beta$$, we obtain the following Double Angle' Identities.
Note Double Angle Identities
For all applicable angles $$\theta$$, \index{Double Angle Identities}
• $$\cos(2\theta) = \left\{ \begin{array}{l} \cos^{2}(\theta) - \sin^{2}(\theta)\\ [5pt] 2\cos^{2}(\theta) - 1 \\ [5pt] 1-2\sin^{2}(\theta) \end{array} \right.$$
• $$\sin(2\theta) = 2\sin(\theta)\cos(\theta)$$
• $$\tan(2\theta) = \dfrac{2\tan(\theta)}{1 - \tan^{2}(\theta)}$$
The three different forms for $$\cos(2\theta)$$ can be explained by our ability to exchange' squares of cosine and sine via the Pythagorean Identity $$\cos^{2}(\theta) + \sin^{2}(\theta) = 1$$ and we leave the details to the reader. It is interesting to note that to determine the value of $$\cos(2\theta)$$, only \textit{one} piece of information is required: either $$\cos(\theta)$$ or $$\sin(\theta)$$. To determine $$\sin(2\theta)$$, however, it appears that we must know both $$\sin(\theta)$$ and $$\cos(\theta)$$. In the next example, we show how we can find $$\sin(2\theta)$$ knowing just one piece of information, namely $$\tan(\theta)$$.
Example $$\PageIndex{1}$$:
1. Suppose $$P(-3,4)$$ lies on the terminal side of $$\theta$$ when $$\theta$$ is plotted in standard position. Find $$\cos(2\theta)$$ and $$\sin(2\theta)$$ and determine the quadrant in which the terminal side of the angle $$2\theta$$ lies when it is plotted in standard position.
2. If $$\sin(\theta) = x$$ for $$-\frac{\pi}{2} \leq \theta \leq \frac{\pi}{2}$$, find an expression for $$\sin(2\theta)$$ in terms of $$x$$.
3. \label{doubleanglesinewtan} Verify the identity: $$\sin(2\theta) = \dfrac{2\tan(\theta)}{1 + \tan^{2}(\theta)}$$.
4. Express $$\cos(3\theta)$$ as a polynomial in terms of $$\cos(\theta)$$.
Solution
1. Using Theorem \ref{cosinesinecircle} from Section \ref{TheUnitCircle} with $$x = -3$$ and $$y=4$$, we find $$r = \sqrt{x^2+y^2} = 5$$. Hence, $$\cos(\theta) = -\frac{3}{5}$$ and $$\sin(\theta) = \frac{4}{5}$$. Applying Theorem \ref{doubleangle}, we get $$\cos(2\theta) = \cos^{2}(\theta) - \sin^{2}(\theta) = \left(-\frac{3}{5}\right)^2 - \left(\frac{4}{5}\right)^2 = -\frac{7}{25}$$, and $$\sin(2\theta) = 2 \sin(\theta) \cos(\theta) = 2 \left(\frac{4}{5}\right)\left(-\frac{3}{5}\right) = -\frac{24}{25}$$. Since both cosine and sine of $$2\theta$$ are negative, the terminal side of $$2\theta$$, when plotted in standard position, lies in Quadrant III.
2. If your first reaction to $\sin(\theta) = x$' is No it's not, $$\cos(\theta) = x!' then you have indeed learned something, and we take comfort in that. However, context is everything. Here, x' is just a variable - it does not necessarily represent the \(x$$-coordinate of the point on The Unit Circle which lies on the terminal side of $$\theta$$, assuming $$\theta$$ is drawn in standard position. Here, $$x$$ represents the quantity $$\sin(\theta)$$, and what we wish to know is how to express $$\sin(2\theta)$$ in terms of $$x$$. We will see more of this kind of thing in Section \ref{ArcTrig}, and, as usual, this is something we need for Calculus. Since $$\sin(2\theta) = 2 \sin(\theta) \cos(\theta)$$, we need to write $$\cos(\theta)$$ in terms of $$x$$ to finish the problem. We substitute $$x = \sin(\theta)$$ into the Pythagorean Identity, $$\cos^{2}(\theta) + \sin^{2}(\theta) = 1$$, to get $$\cos^{2}(\theta) + x^2 = 1$$, or $$\cos(\theta) = \pm \sqrt{1-x^2}$$. Since $$-\frac{\pi}{2} \leq \theta \leq \frac{\pi}{2}$$, $$\cos(\theta) \geq 0$$, and thus $$\cos(\theta) = \sqrt{1-x^2}$$. Our final answer is $$\sin(2\theta) = 2 \sin(\theta) \cos(\theta) = 2x\sqrt{1-x^2}$$.
We start with the right hand side of the identity and note that $$1 + \tan^{2}(\theta) = \sec^{2}(\theta)$$. From this point, we use the Reciprocal and Quotient Identities to rewrite $$\tan(\theta)$$ and $$\sec(\theta)$$ in terms of $$\cos(\theta)$$ and $$\sin(\theta)$$:
$\begin{array}{rcl} \dfrac{2\tan(\theta)}{1 + \tan^{2}(\theta)} & = & \dfrac{2\tan(\theta)}{\sec^{2}(\theta)}= \dfrac{2 \left( \dfrac{\sin(\theta)}{\cos(\theta)}\right)}{\dfrac{1}{\cos^{2}(\theta)}}= 2\left( \dfrac{\sin(\theta)}{\cos(\theta)}\right) \cos^{2}(\theta) \\ & = & 2\left( \dfrac{\sin(\theta)}{\cancel{\cos(\theta)}}\right) \cancel{\cos(\theta)} \cos(\theta) = 2\sin(\theta) \cos(\theta) = \sin(2\theta) \\ \end{array}$
1. In Theorem \ref{doubleangle}, one of the formulas for $$\cos(2\theta)$$, namely $$\cos(2\theta) = 2\cos^{2}(\theta) - 1$$, expresses $$\cos(2\theta)$$ as a polynomial in terms of $$\cos(\theta)$$. We are now asked to find such an identity for $$\cos(3\theta)$$. Using the sum formula for cosine, we begin with
$\begin{array}{rcl} \cos(3\theta) & = & \cos(2\theta + \theta) \\ & = & \cos(2\theta)\cos(\theta) - \sin(2\theta)\sin(\theta) \\ \end{array}$
Our ultimate goal is to express the right hand side in terms of $$\cos(\theta)$$ only. We substitute $$\cos(2\theta) = 2\cos^{2}(\theta) -1$$ and $$\sin(2\theta) = 2\sin(\theta)\cos(\theta)$$ which yields
$\begin{array}{rcl} \cos(3\theta) & = & \cos(2\theta)\cos(\theta) - \sin(2\theta)\sin(\theta) \\ & = & \left(2\cos^{2}(\theta) - 1\right) \cos(\theta) - \left(2 \sin(\theta) \cos(\theta) \right)\sin(\theta) \\ & = & 2\cos^{3}(\theta)- \cos(\theta) - 2 \sin^2(\theta) \cos(\theta) \\ \end{array}$
Finally, we exchange $$\sin^{2}(\theta)$$ for $$1 - \cos^{2}(\theta)$$ courtesy of the Pythagorean Identity, and get
$\begin{array}{rcl} \cos(3\theta) & = & 2\cos^{3}(\theta)- \cos(\theta) - 2 \sin^2(\theta) \cos(\theta) \\ & = & 2\cos^{3}(\theta)- \cos(\theta) - 2 \left(1 - \cos^{2}(\theta)\right) \cos(\theta) \\ & = & 2\cos^{3}(\theta)- \cos(\theta) - 2\cos(\theta) + 2\cos^{3}(\theta) \\ & = & 4\cos^{3}(\theta)- 3\cos(\theta) \\ \end{array}$
and we are done.
In the last problem in Example \ref{doubleangleex}, we saw how we could rewrite $$\cos(3\theta)$$ as sums of powers of $$\cos(\theta)$$. In Calculus, we have occasion to do the reverse; that is, reduce the power of cosine and sine. Solving the identity $$\cos(2\theta) = 2\cos^{2}(\theta) -1$$ for $$\cos^{2}(\theta)$$ and the identity $$\cos(2\theta) = 1 - 2\sin^{2}(\theta)$$ for $$\sin^{2}(\theta)$$ results in the aptly-named `Power Reduction' formulas below.
Power Reduction Formulas
For all angles $$\theta$$, \index{Power Reduction Formulas}
• $$\cos^{2}(\theta) = \dfrac{1 + \cos(2\theta)}{2}$$
• $$\sin^{2}(\theta) = \dfrac{1 - \cos(2\theta)}{2}$$
Example $$\PageIndex{1}$$:
Rewrite $$\sin^{2}(\theta) \cos^{2}(\theta)$$ as a sum and difference of cosines to the first power.
Solution
We begin with a straightforward application of Theorem \ref{powerreduction}
$\begin{array}{rcl} \sin^{2}(\theta) \cos^{2}(\theta) & = & \left( \dfrac{1 - \cos(2\theta)}{2} \right) \left( \dfrac{1 + \cos(2\theta)}{2} \right) \\ & = & \dfrac{1}{4}\left(1 - \cos^{2}(2\theta)\right) \\ & = & \dfrac{1}{4} - \dfrac{1}{4}\cos^{2}(2\theta) \\ \end{array}$
Next, we apply the power reduction formula to $$\cos^{2}(2\theta)$$ to finish the reduction
$\begin{array}{rcl} \sin^{2}(\theta) \cos^{2}(\theta) & = & \dfrac{1}{4} - \dfrac{1}{4}\cos^{2}(2\theta) \\ & = & \dfrac{1}{4} - \dfrac{1}{4} \left(\dfrac{1 + \cos(2(2\theta))}{2}\right) \\ & = & \dfrac{1}{4} - \dfrac{1}{8} - \dfrac{1}{8}\cos(4\theta) \\ & = & \dfrac{1}{8} - \dfrac{1}{8}\cos(4\theta) \\ \end{array}$
Another application of the Power Reduction Formulas is the Half Angle Formulas. To start, we apply the Power Reduction Formula to $$\cos^{2}\left(\frac{\theta}{2}\right) $\cos^{2}\left(\dfrac{\theta}{2}\right) = \dfrac{1 + \cos\left(2 \left(\frac{\theta}{2}\right)\right)}{2} = \dfrac{1 + \cos(\theta)}{2}.$ We can obtain a formula for \(\cos\left(\frac{\theta}{2}\right)$$ by extracting square roots. In a similar fashion, we may obtain a half angle formula for sine, and by using a quotient formula, obtain a half angle formula for tangent. We summarize these formulas below.
Half Angle Formulas
For all applicable angles $$\theta$$:
• $$\cos\left(\dfrac{\theta}{2}\right) = \pm \sqrt{\dfrac{1 + \cos(\theta)}{2}}$$
• $$\sin\left(\dfrac{\theta}{2}\right) = \pm \sqrt{\dfrac{1 - \cos(\theta)}{2}}$$
• $$\tan\left(\dfrac{\theta}{2}\right) = \pm \sqrt{\dfrac{1 - \cos(\theta)}{1+\cos(\theta)}}$$
where the choice of $$\pm$$ depends on the quadrant in which the terminal side of $$\dfrac{\theta}{2}$$ lies.
Example $$\PageIndex{1}$$:
1. Use a half angle formula to find the exact value of $$\cos\left(15^{\circ}\right)$$.
2. Suppose $$-\pi \leq \theta \leq 0$$ with $$\cos(\theta) = -\frac{3}{5}$$. Find $$\sin\left(\frac{\theta}{2}\right)$$.
3. Use the identity given in number \ref{doubleanglesinewtan} of Example \ref{doubleangleex} to derive the identity $\tan\left(\dfrac{\theta}{2}\right) = \dfrac{\sin(\theta)}{1+\cos(\theta)}$
Solution
1. To use the half angle formula, we note that $$15^{\circ} = \frac{30^{\circ}}{2}$$ and since $$15^{\circ}$$ is a Quadrant I angle, its cosine is positive. Thus we have
$\begin{array}{rcl} \cos\left(15^{\circ}\right) & = & + \sqrt{\dfrac{1+\cos\left(30^{\circ}\right)}{2}} = \sqrt{\dfrac{1+\frac{\sqrt{3}}{2}}{2}}\\ & = & \sqrt{\dfrac{1+\frac{\sqrt{3}}{2}}{2}\cdot \dfrac{2}{2}} = \sqrt{\dfrac{2+\sqrt{3}}{4}} = \dfrac{\sqrt{2+\sqrt{3}}}{2}\\ \end{array}$
Back in Example \ref{cosinesumdiffex}, we found $$\cos\left(15^{\circ}\right)$$ by using the difference formula for cosine. In that case, we determined $$\cos\left(15^{\circ}\right) = \frac{\sqrt{6}+ \sqrt{2}}{4}$$. The reader is encouraged to prove that these two expressions are equal.
1. If $$-\pi \leq \theta \leq 0$$, then $$-\frac{\pi}{2} \leq \frac{\theta}{2} \leq 0$$, which means $$\sin\left(\frac{\theta}{2}\right) < 0$$. Theorem \ref{halfangle} gives
$\begin{array}{rcl} \sin\left(\dfrac{\theta}{2} \right) & = & -\sqrt{\dfrac{1-\cos\left(\theta \right)}{2}} = -\sqrt{\dfrac{1- \left(-\frac{3}{5}\right)}{2}}\\ & = & -\sqrt{\dfrac{1 + \frac{3}{5}}{2} \cdot \dfrac{5}{5}} = -\sqrt{\dfrac{8}{10}} = -\dfrac{2\sqrt{5}}{5}\\ \end{array}$
1. Instead of our usual approach to verifying identities, namely starting with one side of the equation and trying to transform it into the other, we will start with the identity we proved in number \ref{doubleanglesinewtan} of Example \ref{doubleangleex} and manipulate it into the identity we are asked to prove. The identity we are asked to start with is $$\; \sin(2\theta) = \frac{2\tan(\theta)}{1 + \tan^{2}(\theta)}$$. If we are to use this to derive an identity for $$\tan\left(\frac{\theta}{2}\right)$$, it seems reasonable to proceed by replacing each occurrence of $$\theta$$ with $$\frac{\theta}{2} $\begin{array}{rcl} \sin\left(2 \left(\frac{\theta}{2}\right)\right) & = & \dfrac{2\tan\left(\frac{\theta}{2}\right)}{1 + \tan^{2}\left(\frac{\theta}{2}\right)} \\ \sin(\theta) & = & \dfrac{2\tan\left(\frac{\theta}{2}\right)}{1 + \tan^{2}\left(\frac{\theta}{2}\right)} \\ \end{array}$ We now have the \(\sin(\theta)$$ we need, but we somehow need to get a factor of $$1+\cos(\theta)$$ involved. To get cosines involved, recall that $$1 + \tan^{2}\left(\frac{\theta}{2}\right) = \sec^{2}\left(\frac{\theta}{2}\right)$$. We continue to manipulate our given identity by converting secants to cosines and using a power reduction formula
$\begin{array}{rcl} \sin(\theta) & = & \dfrac{2\tan\left(\frac{\theta}{2}\right)}{1 + \tan^{2}\left(\frac{\theta}{2}\right)} \\ \sin(\theta) & = & \dfrac{2\tan\left(\frac{\theta}{2}\right)}{\sec^{2}\left(\frac{\theta}{2}\right)} \\ \sin(\theta) & = & 2 \tan\left(\frac{\theta}{2}\right) \cos^{2}\left(\frac{\theta}{2}\right) \\ \sin(\theta) & = & 2 \tan\left(\frac{\theta}{2}\right) \left(\dfrac{1 + \cos\left(2 \left(\frac{\theta}{2}\right)\right)}{2}\right) \\ \sin(\theta) & = & \tan\left(\frac{\theta}{2}\right) \left(1+\cos(\theta) \right) \\ \tan\left(\dfrac{\theta}{2}\right) & = & \dfrac{\sin(\theta)}{1+\cos(\theta)} \\ \end{array}$
Our next batch of identities, the Product to Sum Formulas,\footnote{These are also known as the Prosthaphaeresis Formulas and have a rich history. The authors recommend that you conduct some research on them as your schedule allows.} are easily verified by expanding each of the right hand sides in accordance with Theorem \ref{circularsumdifference} and as you should expect by now we leave the details as exercises. They are of particular use in Calculus, and we list them here for reference.
Note: Product to Sum Formulas
For all angles $$\alpha$$ and $$\beta$$, \index{Product to Sum Formulas}
• $$\cos(\alpha)\cos(\beta) = \frac{1}{2} \left[ \cos(\alpha - \beta) + \cos(\alpha + \beta)\right]$$
• $$\sin(\alpha)\sin(\beta) = \frac{1}{2} \left[ \cos(\alpha - \beta) - \cos(\alpha + \beta)\right]$$
• $$\sin(\alpha)\cos(\beta) = \frac{1}{2} \left[ \sin(\alpha - \beta) + \sin(\alpha + \beta)\right]$$
Related to the Product to Sum Formulas are the Sum to Product Formulas, which we will have need of in Section \ref{TrigEquIneq}. These are easily verified using the Product to Sum Formulas, and as such, their proofs are left as exercises.
Note: Sum to Product Formulas:
For all angles $$\alpha$$ and $$\beta$$:
1. $$\cos(\alpha) + \cos(\beta) = 2 \cos\left( \dfrac{\alpha + \beta}{2}\right)\cos\left( \dfrac{\alpha - \beta}{2}\right)$$
2. $$\cos(\alpha) - \cos(\beta) = - 2 \sin\left( \dfrac{\alpha + \beta}{2}\right)\sin\left( \dfrac{\alpha - \beta}{2}\right)$$
3. $$\sin(\alpha) \pm \sin(\beta) = 2 \sin\left( \dfrac{\alpha \pm \beta}{2}\right)\cos\left( \dfrac{\alpha \mp \beta}{2}\right)$$
Example $$\PageIndex{1}$$:
1. Write $$\; \cos(2\theta)\cos(6\theta) \;$$ as a sum.
2. \Write $$\; \sin(\theta) - \sin(3\theta) \;$$ as a product.
Solution
1. Identifying $$\alpha = 2\theta$$ and $$\beta = 6\theta$$, we find
$\begin{array}{rcl} \cos(2\theta)\cos(6\theta) & = & \frac{1}{2} \left[ \cos(2\theta - 6\theta) + \cos(2\theta + 6\theta)\right]\\ & = & \frac{1}{2} \cos(-4\theta) + \frac{1}{2}\cos(8\theta) \\ & = & \frac{1}{2} \cos(4\theta) + \frac{1}{2} \cos(8\theta), \end{array}$
where the last equality is courtesy of the even identity for cosine, $$\cos(-4\theta) = \cos(4\theta)$$.
1. Identifying $$\alpha = \theta$$ and $$\beta = 3\theta$$ yields
$\begin{array}{rcl} \sin(\theta) - \sin(3\theta) & = & 2 \sin\left( \dfrac{\theta - 3\theta}{2}\right)\cos\left( \dfrac{\theta + 3\theta}{2}\right) \\ & = & 2 \sin\left( -\theta \right)\cos\left( 2\theta \right) \\ & = & -2 \sin\left( \theta \right)\cos\left( 2\theta \right), \\ \end{array}$
where the last equality is courtesy of the odd identity for sine, $$\sin(-\theta) = -\sin(\theta)$$.
The reader is reminded that all of the identities presented in this section which regard the circular functions as functions of angles (in radian measure) apply equally well to the circular (trigonometric) functions regarded as functions of real numbers. In Exercises \ref{idengraphfirst} - \ref{idengraphlast} in Section \ref{TrigGraphs}, we see how some of these identities manifest themselves geometrically as we study the graphs of the these functions. In the upcoming Exercises, however, you need to do all of your work analytically without graphs.
## Contributors
• Carl Stitz, Ph.D. (Lakeland Community College) and Jeff Zeager, Ph.D. (Lorain County Community College) | 2020-01-27T18:17:14 | {
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https://stats.stackexchange.com/questions/336321/when-is-the-law-of-product-of-probabilities-applicable | # When is the law of product of probabilities applicable?
The sample space of the experiment of throwing a green and a red dice has 36 elements. The event, say $A$, that the sum $x+y>8$ in $(x,y)$, where $x$ is an outcome of the green die and $y$ that of red die will occur has the probability:
$$P(A)=\frac{10}{36}=\frac{5}{18}$$
And the probability of the event, say C, that a number greater than 4 will turn up on the green die is:
$$P(C)=\frac{12}{36} = \frac{1}{3}$$
And the probability of the intersection $A \cap C$ is:
$$P(A \cap C) = \frac{7}{36}$$
which is not equal to the product of the probabilities of $A$ and $C$, that is
$$P(A) \cdot P(C)=\frac{5}{54} \neq \frac{7}{36}$$
Is the law of product even applicable here? If so, how?
• The problem is that the events A and C that you selected are NOT independent, So the product rule may not apply. – Michael Chernick Mar 23 '18 at 17:01
The events A & C are dependent event. That is, if one of them occurs first, the odds of second event are changed. The law of product for dependent events is:
P(A∩C) = P(C).P(A|C)
Where P(A|C) is conditional probability of the event A given event C has already occurred.
If event C (x > 4) has already occurred, then we have a sample space of 12 instead of 36. This is already clear in your calculation of P(C).
There are 7 Dice throw cases where event A (x + y > 8) occurs, given that event C (x > 4) already occurred: (5,4) (5,5) (5,6) (6,3) (6,4) (6,5) (6,6)
hence P(A|C) = 7/12
P(A∩C) = P(C).P(A|C) = 1/3 . 7/12 = 7/36
which is same as the P(A∩C) calculated independently.
By the way, we can also prove the same other way around, that is: P(A∩C) = P(A).P(C|A) (But I chose the case which I was comfortable with.)
PS: pls disregard my clumsy writing, i still have not adapted to write fractions & other mathematic symbols in this forum.
Product of probabilities equals their joint probability only for independent events. In fact, it is a part of definition of independence:
Two events A and B are independent (often written as $A \perp B$ or $A \perp\!\!\!\perp B$) if their joint probability equals the product of their probabilities:
$$\mathrm{P}(A \cap B) = \mathrm{P}(A)\mathrm{P}(B)$$
• Could you explain why $A$ and $C$ are not independent events? – Samama Fahim Mar 23 '18 at 17:44
• @SamamaFahim because the first event "tells" you something about possible outcomes of the first one. Imagine more extreme example A=z<5 and B=z<4.9999..., if I told you that A is true, would it change anything about your bet on B? – Tim Mar 23 '18 at 17:57 | 2019-08-21T03:19:43 | {
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https://math.stackexchange.com/questions/2573388/expected-value-of-position-in-subset | # Expected value of position in subset.
I have set of $n$ elements $[0, ... n-1]$.
I randomly pick a subset $S$ of $k$ elements (also ordered).
Assume I have $t \in \{1, .. k\}$. What is expected value of $t$-th position in ordered subset?
I try to solve it like this:
$E\big[S_{t}\big] = \sum_{i=1}^{n} \big((i - 1) * P(S_{t}=(i-1))\big)$
$P(S_{t} = i) = \cfrac{\binom{i - 1}{t-1} * \binom{n-i}{k-t}}{\binom{n}{k}} .$
Here:
$\binom{i - 1}{t-1}$ is for selecting prefix;
$\binom{n-i}{k-t}$ is for postfix;
But I can't simplify this equality much further (I guess that sum of product of combinations could be simplified).
• I think you will find it easier if you start indexing from $1$; i.e. considering the set $\{1,\ldots, n\}$, and also assuming $t \in \{1,\ldots, k\}$: this will facilitate working with Binomials. Can you also explain why you think you need the factor $\frac1n$ in the probability. Have you tried some explicit examples? Eg. n = 5, k = 3, and t = 2? – owen88 Dec 19 '17 at 17:47
• It was my last thought about $\frac{1}{n}$, may be it's not valid (thought it's like select an element $i$ and then select prefix and postfix). – Sindbag Dec 19 '17 at 17:55
• You do not need the `select an element i' part, as this is assumed in the statement $P[S_t = i]$; i.e. you are counting those events for which the $t$-th element is $i$; otherwise your logic is correct. It'll help if you check some small examples as I prompted above. – owen88 Dec 19 '17 at 18:44
As in my comment above, I will consider the set to be $\{1,\ldots, n\}$, and let $t \in \{1,\ldots, k\}$ denote the position of interest, i.e. $S_t$ is the value of the $t$-th element of the random subset $S$, when in increasing order.
In this case, we can show that
$$P(S_t = i) = \frac{\binom{i-1}{t-1} \binom{n-i}{k-t} }{\binom{n}{k} }$$
Now substituting this into the formula for expectation:
\begin{align*} E[S_t] & = \sum_{i=1}^n i \frac{\binom{i-1}{t-1} \binom{n-i}{k-t} }{\binom{n}{k} } \\ & = \frac{1}{\binom{n}{k} } \sum_{i = 1}^n i \binom{i-1}{t-1} \binom{n-i}{k-t} \end{align*}
Using the formula: $$i \binom{i-1}{t-1} = t \binom{i}{t}$$
then the summation can be written as:
\begin{align*} \sum_{i=1}^n i \binom{i-1}{t-1}\binom{n-i}{k-t} &= t \sum_{i=1}^n \binom{i}{t}\binom{n-i}{k-t} \\ & = t \left( \sum_{i=0}^n \binom{i}{t}\binom{n-i}{k-t} - \binom{0}{t}\binom{n}{k-t} \right) \end{align*} where we have extended the sum over the range $i = 0,\ldots, n$ (instead of $i = 1, \ldots, n$). The first of the two terms is identified as the Vandermonde identity, whilst the second simplifies to $0$ so that
\begin{align*} \sum_{i=1}^n i \binom{i-1}{t-1}\binom{n-i}{k-t} &= t \binom{n+1}{k+1} \end{align*}
So over all, we have:
\begin{align*} E[S_t] &= t \frac{ \binom{n+1}{k+1} }{\binom{n}{k}} \\ & = t \frac{n+1}{k+1} \end{align*}
• Are you sure about Vandermonde's identity? We have $i$ changing, whereas Vandermonde's requires $t$ to be changing. – Sindbag Dec 19 '17 at 18:56
• See equation (9), en.wikipedia.org/wiki/… If you agree with my answer, will you please like and accept it. Thanks – owen88 Dec 19 '17 at 19:07
• Great, thanks a lot! Can't like it - too little rating :( – Sindbag Dec 19 '17 at 19:17
• No problem! Glad to have helped. – owen88 Dec 19 '17 at 19:18 | 2019-05-23T14:56:42 | {
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https://math.stackexchange.com/questions/2968435/why-is-an-element-in-mathbbz-sqrt2-a-unit-iff-a2-2b2-pm-1/2968442 | # Why is an element in $\mathbb{Z}[\sqrt{2}]$ a unit iff $a^2-2b^2 = \pm 1$?
In this answer and other's that I've read, finding the units in $$\mathbb{Z}[\sqrt{2}]$$ (elements of the form $$a + b\sqrt{2}$$ for integers $$a$$ and $$b$$) is equivalent to finding the solutions to $$a^2-2b^2 = \pm 1$$. So from my understanding, for a ring $$R$$, $$u \in R$$ is a unit if there exists a $$v \in R$$ such that $$uv = vu = 1_R$$, with $$1_R$$ being the multiplicative identity in the ring.
So I have 2 questions:
1. My understanding is that we arrive at this because $$(a+b\sqrt{2})(a-b\sqrt{2}) = a^2-2b^2$$, so units can only come in conjugate pairs. I could intuitively reason that this is because $$(a+b\sqrt{2})x$$ would have a $$\sqrt{2}$$ term if $$x$$ is not conjugate, but I'm not sure if this is strictly the reason why the units are only conjugates.
2. Why can they multiply to $$-1$$ and still be a unit? Is it because if $$(a+b\sqrt{2})(a-b\sqrt{2}) = a^2-2b^2 = -1$$, then $$(a+b\sqrt{2})^2(a-b\sqrt{2})^2 = (-1)^2 = 1$$, and thus the inverse of $$(a+b\sqrt{2})$$ is $$(a+b\sqrt{2})(a-b\sqrt{2})^2$$?
• Are you familiar with the norm $N(x)$ defined on $\mathbb Z[\sqrt 2]$? – Frpzzd Oct 23 '18 at 23:31
• @Frpzzd I am not - this question came up in one of my first lectures on the topic of rings, so I'm very new to it. – rb612 Oct 23 '18 at 23:32
• Okay. I will define it for you in my answer. – Frpzzd Oct 23 '18 at 23:33
• Have I explained by answer sufficiently? Please let me know if anything in my answer is confusing or worded badly so that I can clarify. – Frpzzd Oct 23 '18 at 23:43
Define the norm $$N:\mathbb Z[\sqrt 2]\mapsto \mathbb Z$$ as follows: $$N(a+b\sqrt 2)=a^2-2b^2\tag{1}$$ or $$N(x)=x\bar{x}\tag{2}$$ Then notice that for all $$x,y\in\mathbb Z[\sqrt 2]$$, we have that $$N(xy)=N(x)N(y)\tag{3}$$ Thus, suppose that $$x\in\mathbb Z[\sqrt 2]$$ is a unit. Then let $$x^{-1}$$ be its inverse. We have by $$(3)$$ that $$1=N(1)=N(xx^{-1})=N(x)N(x^{-1})$$ and so $$N(x)$$, which is an integer, is a unit, with multiplicative inverse $$N(x^{-1})$$. Since the only units in $$\mathbb Z$$ are $$\pm 1$$, we have that $$N(x)=a^2-2b^2$$ must equal $$+1$$ or $$-1$$.
This proves that $$a^2-2b^2$$ must equal $$1$$ or $$-1$$ if $$x$$ is a unit. Now we must prove the converse.
Suppose $$N(x)=a^2-2b^2=\pm 1$$. Then, by $$(2)$$, $$x\bar{x}=1$$ or $$x\bar{x}=-1$$. If the former is true, then $$\bar{x}$$ is the inverse of $$x$$ and $$x$$ is a unit; if the latter is true, then $$-\bar{x}$$ is the inverse of $$x$$ and $$x$$ is a unit.
Proof complete!
Now I shall address a generalized form of your first question. Suppose that $$x,y\in\mathbb Z[\sqrt 2]$$, and that $$xy\in\mathbb Z$$. Write $$x=a+b\sqrt 2$$ and $$y=c+d\sqrt 2$$. Then $$xy=(ac+2bd)+(ad+bc)\sqrt{2}$$ For $$xy\in\mathbb Z$$ to hold, it must be true that the coefficient $$ad+bc$$ of $$\sqrt 2$$ vanishes. Thus, we have $$ad+bc=0$$ or $$\frac{a}{b}=-\frac{c}{d}$$ ...which implies that the ratios of the rational and irrational parts of $$x$$ and $$y$$ are equal and opposite. This is equivalent to stating the existence of $$z\in\mathbb Z[\sqrt 2]$$ such that $$z|x$$ and $$\bar{z}|y$$.
If you apply this to units, you can see that it is indeed true that all units must come in conjugate pairs.
• Thanks so much for the thorough response! This is great. My main question from this though is why can we/do define the norm as such and then proceed from there? In other words, did we arbitrarily define that norm to proceed with the proof, or is it the only option for the norm of an element in this ring? – rb612 Oct 23 '18 at 23:51
• @rb612 We choose that norm mostly because of its nice multiplicative property $N(xy)=N(x)N(y)$. It's not really that arbitrary; I've just assigned the value $a^2-2b^2$ mentioned in your question to a function $N(x)$, giving it a name. I could have done the proof without defining $N(x)$ at all and instead referring to $a^2-2b^2$; it's just the convention that I am used to using. – Frpzzd Oct 23 '18 at 23:53
• @rb612 One more thing... once I've answered all of your questions, don't forget to $\checkmark$! :D – Frpzzd Oct 23 '18 at 23:55 | 2019-03-19T18:50:41 | {
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https://math.stackexchange.com/questions/3704345/will-the-set-of-algebraic-polynomials-be-dense-in-space-l-infty0-1 | # Will the set of algebraic polynomials be dense in space $L_{\infty}(0,1)$?
I know that the set of algebraic polynomials is dense in space $$L_{p}(a,b)$$, where $$1 \leq p < \infty$$ and $$a,b \in \mathbb{R}$$. However, what about $$L_{\infty}(0,1)$$? Will the set of algebraic polynomials be dense in space $$L_{\infty}(0,1)$$?
In my opinion it is not, but is there any example that prove this fact?
• Since $L_{\infty}(0,1)$ is not separable, it cannot have the set of real polynomials as a dense subset. If it were true, the set of algebraic polynomials with rational coefficients would be dense in $L_{\infty}(0,1)$. But that is not possible since the set of rational polynomials is countable. – Darth Lubinus Jun 3 at 18:48
• @DarthLubinus, Could you explain me why this is impossible? – Sh VavilenT Jun 3 at 19:52
• Sure, I'll give some details in an answer! – Darth Lubinus Jun 3 at 19:55
I'm going to try and ellaborate on my previous comment.
First, let's call $$\mathbb{R}[X]$$ the ring of real polynomials, and suppose $$\mathbb{R}[X]$$ was a dense subset of $$L_{\infty}(0,1)$$. Therefore, we can prove that $$\mathbb{Q}[X]$$ (the subset of all rational polynomials) is dense in $$L_{\infty}(0,1)$$ by seeing that it is dense in $$\mathbb{R}[X]$$.
Let $$p(x)=a_{n}x^{n}+...+a_{1}x+a_{0}$$ be a polynomial with coefficients in $$\mathbb{R}$$, and $$\varepsilon>0$$. Then, for every $$a_{i}$$ we can find a $$b_{i}\in \mathbb{Q}$$ such that $$|b_{i}-a_{i}|<\varepsilon/n$$. By taking $$q(x)=b_{n}x^{n}+...+b_{1}x+b_{0}$$, we have that $$q\in \mathbb{Q}[X]$$, and for every $$x\in [0,1]$$,
$$|p(x)-q(x)|=\Big|\sum_{i=0}^{n}(a_{i}-b_{i})x\Big|\leq \sum_{i=0}^{n}|a_{i}-b_{i}||x|\leq n\frac{\varepsilon}{n}=\varepsilon,$$
so $$||p-q||_{\infty}\leq\varepsilon$$. Therefore, $$\mathbb{Q}[X]$$ is dense in $$\mathbb{R}[X]$$ and, using our assumption that $$\mathbb{R}[X]$$ is dense in $$L_{\infty}(0,1)$$, we conclude that $$\mathbb{Q}[X]$$ is dense in $$L_{\infty}(0,1)$$. In particular, $$L_{\infty}(0,1)$$ admits a countable dense subspace, which means that it is separable.
To find our desired contradiction, we need to prove that $$L_{\infty}(0,1)$$ cannot be separable.
Consider the family $$\{f_{t}\colon t\in (0,1)\}\subseteq L_{\infty}(0,1)$$, where $$f_{t}(s)=1$$ if $$s\leq t$$, and $$f_{t}(s)=0$$ otherwise. It is easy to see that, whenever $$t\neq s$$, we have $$||f_{t}-f_{s}||_{\infty}=1$$. Therefore, we can find an uncountable family of disjoint open balls $$A_{t}=B(f_{t},1/3)$$, where $$0 < t < 1$$. Since $$\mathbb{Q}[X]$$ is dense in $$L_{\infty}(0,1)$$, each $$A_{t}$$ must contain at least one polynomial $$q_{t}\in \mathbb{Q}[X]$$, but this is impossible, because it would imply that there is an uncountable amount of rational polynomials! Therefore, $$\mathbb{R}[X]$$ cannot be dense.
In essence, what is happening is that $$L_{\infty}$$ is too big of a space to be the closure of something as comparably small as a countable subset. When $$p<\infty$$, this doesn't happen, and the fact that polynomials are dense in $$L_{p}(0,1)$$ proves that this last space is separable.
Just for the sake of adding another point of view, here is a more "topological" proof (not too different from the one I just did). In metric spaces, separability is the same thing as second countability (the property of having a countable basis for the topology of the metric space). If $$L_{\infty}(0,1)$$ were separable/secound countable, then the family $$B=\{f_{t}\}_{t\in (0,1)}$$ would be another second countable space (as a subspace of one). But $$||f_{t}-f_{s}||_{\infty}=1$$ when $$t\neq s$$, so $$B$$ must be discrete. In that case, $$B$$ must be countable, which is again a contradiction.
I hope this helps! | 2020-08-10T22:41:05 | {
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http://mathhelpforum.com/calculus/256026-volume-solid-revolution.html | # Thread: Volume of the solid of revolution
1. ## Volume of the solid of revolution
There are a region bounded by the curve y=cos^2(x)sin^4(x) and the x-axis , where pi<x<2pi. Find the volumn of the solid of revolution when the region is revolved about the y-axis
2. ## Re: Volume of the solid of revolution
Hi Lau, what have you tried already? It's not clear if you are stuck on what method to use to solve this or if the integration is too difficult.
3. ## Re: Volume of the solid of revolution
Hey Lau, I solved this problem using the shell method.
It is best to start with a rough sketch of the function on the given interval (Will give a M shape), and once you do you will realize that we should split the region again this time considering: pi <x<3pi/2 and 3pi/2<x<2pi. This is because from x=pi to x=3pi/2 , the function "starts" then "ends". Also, this will help us in the calculations as the function "repeats" its self on 3pi/2 < x < 2pi.
Using the shell method, I found the height of a shell to be: $\displaystyle h=f(x)= \sin^4x \cos^2x$ and a radius about the y axis to be $\displaystyle r=x$.
Then the volume of the first cycle of the periodic function is:
$\displaystyle V= \int_{\pi}^{3 \pi/2} 2 \pi r \cdot h \, dx$
$\displaystyle V= \int_{\pi}^{3 \pi/2} 2 \pi (x) \cdot (\sin^4x \cos^2x) \, dx = \frac{\pi}{36} + \frac{5 \pi}{64}$
So finally, the total volume is just $\displaystyle 2V = \frac{\pi}{18} + \frac{5 \pi}{32} \approx 5.02$
4. ## Re: Volume of the solid of revolution
Originally Posted by sakonpure6
Hey Lau, I solved this problem using the shell method.
It is best to start with a rough sketch of the function on the given interval (Will give a M shape), and once you do you will realize that we should split the region again this time considering: pi <x<3pi/2 and 3pi/2<x<2pi. This is because from x=pi to x=3pi/2 , the function "starts" then "ends". Also, this will help us in the calculations as the function "repeats" its self on 3pi/2 < x < 2pi.
Using the shell method, I found the height of a shell to be: $\displaystyle h=f(x)= \sin^4x \cos^2x$ and a radius about the y axis to be $\displaystyle r=x$.
Then the volume of the first cycle of the periodic function is:
$\displaystyle V= \int_{\pi}^{3 \pi/2} 2 \pi r \cdot h \, dx$
$\displaystyle V= \int_{\pi}^{3 \pi/2} 2 \pi (x) \cdot (\sin^4x \cos^2x) \, dx = \frac{\pi}{36} + \frac{5 \pi}{64}$
So finally, the total volume is just $\displaystyle \color{red}{2V = \frac{\pi}{18} + \frac{5 \pi}{32} \approx 5.02}$
@sakonpure6, you have a conceptual mistake.
$\displaystyle 3.30405 \approx \int_{1.5\pi }^{2\pi } {\left( {2\pi x} \right)\left( {{{\sin }^4}(x){{\cos }^2}(x)} \right)dx} \ne \int_\pi ^{1.5\pi } {\left( {2\pi x} \right)\left( {{{\sin }^4}(x){{\cos }^2}(x)} \right)dx \approx 2.0963}$
Look at the graph again.
5. ## Re: Volume of the solid of revolution
Originally Posted by Plato
@sakonpure6, you have a conceptual mistake.
$\displaystyle 3.30405 \approx \int_{1.5\pi }^{2\pi } {\left( {2\pi x} \right)\left( {{{\sin }^4}(x){{\cos }^2}(x)} \right)dx} \ne \int_\pi ^{1.5\pi } {\left( {2\pi x} \right)\left( {{{\sin }^4}(x){{\cos }^2}(x)} \right)dx \approx 2.0963}$
Look at the graph again.
EDIT a Typo: 2.50963 not 2.0963
6. ## Re: Volume of the solid of revolution
Oh thanks for pointing that out Plato! So, Lau the correct answer is: 3.304 +2.096 = 5.40
Edit: Plato, did something in the equation give away that the function was not periodic - had different amplitudes? Or did you plot to verify?
7. ## Re: Volume of the solid of revolution
It is best to start with a rough sketch of the function on the given interval (Will give a M shape), and once you do you will realize that we should split the region again this time considering: pi <x<3pi/2 and 3pi/2<x<2pi. This is because from x=pi to x=3pi/2 , the function "starts" then "ends". Also, this will help us in the calculations as the function "repeats" its self on 3pi/2 < x < 2pi.
8. ## Re: Volume of the solid of revolution
Originally Posted by dungtuyet
It is best to start with a rough sketch of the function on the given interval (Will give a M shape), and once you do you will realize that we should split the region again this time considering: pi <x<3pi/2 and 3pi/2<x<2pi. This is because from x=pi to x=3pi/2 , the function "starts" then "ends". Also, this will help us in the calculations as the function "repeats" its self on 3pi/2 < x < 2pi.
This is wrong. In fact, it completely misses the point of this whole thread!
@dungtuyet, did you even read the entire thread? If you did, what part did you not get?
One more time. Here is a standard example.
What is the volume of the graph of $y=|x-3|,~ 2\le x\le 4$ about the $y \text{-axis}~?$
Please study each of these: from $x=2\to 3$; then from $x=3\to 4$
Note that those two volumes are not equal even though the graph is symmetric about $x=3$. The radii are different?
Add those two values and compare the result with this. | 2019-07-22T20:39:00 | {
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https://math.stackexchange.com/questions/1186124/how-to-find-the-general-solution-to-this-ode/1186134 | # How to find the general solution to this ODE…
Q:
Find the general solution to the ODE:$$t^2y''-ty'+y=0~~~~~~~~(*)$$given that $y_1=t~~$is a solution.
My intuitive solution:
Let $x=ln(t)$, then $y'=\frac{dy}{dx}\frac{1}{t}$ and $y''=\frac{d^2y}{dx^2}\frac{1}{t^2}-\frac{dy}{dx}\frac{1}{t^2}$
Plug into $(*)$ we get: $$t^2(\frac{d^2y}{dx^2}\frac{1}{t^2}-\frac{dy}{dx}\frac{1}{t^2})-t(\frac{dy}{dx}\frac{1}{t})+y=0$$$$\rightarrow~~\frac{d^2y}{dx^2}-2\frac{dy}{dx}+y=0$$ So the characteristic equation for this is $r^2-2r+1=0~\rightarrow~(r-1)^2=0~\rightarrow~r=1$.
Then $y_1(x)=y_2(x)=e^x~\rightarrow~y_1(t)=y_2(t)=e^{ln(t)}=t.~$But now the Wronskian is certainly zero. So I tried to multiply a $t$ such that $y_2(t)=t^2.~$But $y'_2=2t~~$and$~~ y''_2=2.~$Substituting these into $(*)$ I found that this doesn't solve the equation.
Could someone please give me a hint on how to find the general solution?
• Since $r=1$ is a double-root try $y = Cxe^x$ for the second solution (the first solution is as you have found $y=e^x$). – Winther Mar 12 '15 at 0:26
• Oh...OK, got it, I shouldn't have multiplied the factor to $y_2(t)$ but $y_2(x)$... Thanks! – Ruihong Yuan Mar 12 '15 at 0:47
Based on your analysis, as the solutions of the characteristic equation is $r=1$(double root), the two basic solutions for this ODE are $y_1(x)=e^x\rightarrow y_1(t)=t$ and $y_2(x)=xe^x\rightarrow y_2(t)=t\ln|t|$. Substituting these solutions into $(*)$ you will find they are the solutions. And their combinations are also the solutions, i.e., $$y(t)=C_1t+C_2t\ln|t|$$ where $C_1$ and $C_2$ are constants.
HINT
Consider the method of the reduction of order.
You know one of the solutions to the differential equation, which is: $y_1 = t$
Now, call the second solution $y_2$ where $y_2 = vy_1 = vt$, where $v$ is a function of $t$.
Taking the derivative, we get:
$$y_2 = vt$$
$$y_2' = v't + v$$
$$y_2'' = v''t + v' + v' = v''t + 2v'$$
Substituting this into our equation:
$$t^2(v''t + 2v') -t(v't + v) + vt = 0$$
$$t^3v'' + 2t^2v' - t^2v' - vt + vt =0,\ t^3v'' + 2t^2v' - t^2v' = 0$$
Note: Upon using this method the $y$ term has dropped out. This method ensures that, and now we can find our second solution:
$$t^3v'' + t^2v' = 0$$
Now, we let $w = v'$, $w' = v''$
It follows that:
$$t^3w' + t^2w = 0$$
$$w' = -\frac{w}{t}$$
$$\frac{dw}{w} = -\frac{dt}{t}$$
$$\ln w = -\ln t$$
$$w = -t,\ v' = -t,\ v = -\frac{t^2}{2}$$
Now, our second solution is $$y_2 = vt = -\frac{t^3}{2}$$
Now, our general solution is formed:
$$y(t) = c_1t + c_2\frac{t^3}{2}$$
• From $\ln w=-\ln t$ it follows $w=1/t$, you went a little fast there. Then $v=\ln|t|$ in accordance with the other answers. – Lutz Lehmann Mar 12 '15 at 3:06
• @LutzL ty for catching my mistake – Varun Iyer Mar 12 '15 at 23:23
this is an example of eulers equation. try $y = t^k$ sub it in the equation. find $$k(k-1) - k + 1= 0\to k = 1, 1.$$ the two solutions are $$t, t \ln t$$ and the general solution is $$y = At + B t\ln t.$$
• You meant the type is "Euler equation". – Lutz Lehmann Mar 12 '15 at 3:07 | 2020-02-21T19:42:58 | {
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https://stats.stackexchange.com/questions/103800/calculate-probability-area-under-the-overlapping-area-of-two-normal-distributi/103808 | # Calculate probability (area) under the overlapping area of two normal distributions
I have two normal distributions defined by their averages and standard deviations.
Sample 1: Mean=5.28; SD=0.91
Sample 2: Mean=8.45; SD=1.36
You can see how they look like in the next image:
How can I get the probability to obtain an individual from the overlapping area (green)? Is the probability the same as the area?
• What do you mean by the probability of obtaining an individual from the area? – Juho Kokkala Jun 18 '14 at 8:37
• If you sample points from either normal distribution, you get points on the Perikymata-axis rather than on the 2-dimensional area. Furthermore, the green zone is infinitely wide, so all values sampled from either distribution are under the green zone, so in this sense the probability would be 1. – Juho Kokkala Jun 18 '14 at 8:47
• @JuhoKokkala OP likely means integral over the green area. – Vladislavs Dovgalecs Nov 21 '17 at 18:40
• @VladislavsDovgalecs What would the question "Is the probability the same as the area" then mean? – Juho Kokkala Nov 21 '17 at 19:21
• @VladislavsDovgalecs I think you misunderstood my comment, I was not asking for an interpretation of the area but trying to parse the question. In any case, please note OP's comments to the accepted answer and a followup question posted here stats.stackexchange.com/questions/103821. – Juho Kokkala Nov 23 '17 at 6:14
It is not quite clear what you mean by probability to obtain an individual from the overlapping area. This solves for the area of the green zone in your diagram:
Let:
• $X_1 \sim N(\mu_1,\sigma_1^2)$ with pdf $f_1(x_1)$ and cdf $F_1(x_1)$ and
• $X_2 \sim N(\mu_2,\sigma_2^2)$ with pdf $f_2(x_2)$ and cdf $F_2(x_2)$,
where $\mu_1 < \mu_2$. In your example, the 'black variable' corresponds to $X_1$.
Let $c$ denote the point of intersection where the pdf's meet in the green zone of your plot Then, the area of your green intersection zone is simply:
$$P(X_1>c) + P(X_2<c) = 1 - F_1(c) + F_2(c) = 1-\frac{1}{2} \text{erf}\left(\frac{c-\mu _1}{\sqrt{2} \sigma _1}\right)+\frac{1}{2} \text{erf}\left(\frac{c-\mu _2}{\sqrt{2} \sigma _2}\right)$$
where erf(.) is the error function.
Point $c$ is the solution to $f_1(x) = f_2(x)$ within the green zone, which yields:
$$c = \frac{\mu _2 \sigma _1^2-\sigma _2 \left(\mu _1 \sigma _2+\sigma _1 \sqrt{\left(\mu _1-\mu _2\right){}^2+2 \left(\sigma _1^2-\sigma _2^2\right) \log \left(\frac{\sigma _1}{\sigma _2}\right)}\right)}{\sigma _1^2-\sigma _2^2}$$
For your example, with ${\mu_1 = 5.28, \mu_2 = 8.45, \sigma_1 = 0.91, \sigma_2 = 1.36}$, this yields: $c = 6.70458...$,
and the area of the green section is: 0.158413 ...
• @antecessor Depends on what you mean by 'overlapping among both'. In the question you talk about some probability, but it is unclear what probabilistic interpretation the area computed here would have. – Juho Kokkala Jun 18 '14 at 8:57
• @JuhoKokkala Ok, I will write a new question to address the exact question. – antecessor Jun 18 '14 at 10:38
• the resolution of the equation pdf1=pdf2 yields to two solutions, so we have two intersection points, why you delete the other point @wolfies – user61828 Dec 1 '14 at 11:51
• Because, within the green zone, there is only one solution. – wolfies Dec 1 '14 at 12:32
• Is the equation for c given at the bottom of the answer correct? If both distributions have equal variance, the equation requires dividing by zero! – sammosummo Mar 21 '18 at 21:28
@abdelbasset, To improve @wolfies 's answer above, there are two intersection points c, let's call them c1 and c2. Here, c1=-1.2848 and c2=6.7046. @wolfies ignored c1 (it likely will be too far in the tails of the probability density functions to matter (especially if we round to just a few decimal points)).
The more correct way is to find both c1 and c2, and to find the area of overlap of both functions:
AREA OF OVERLAP = P(X1 > j1) + P(X2 < j1) - [P(X2 < j2)-P(X1 < j2)]
therefore, using cdf's
AREA OF OVERLAP = 1 - F1(j1,μ2,σ2)+F2(j1,μ2,σ2)-F2(j2,μ2,σ2)+F1(j2,μ1,σ1)
• Welcome to the site, @luckapani. You cannot use the "Your Answer" field to comment or respond to comments. As a result, this would be deleted. However, there is the core of an answer here. Can you edit your post to make it more of an answer & not a comment? Since you're new here, you may want to take our tour, which contains information for new users. – gung - Reinstate Monica May 19 '15 at 16:12
Normalise the graphs to an area of 1 by dividing each by their respective standard deviation. Then use simple subtraction from a z-graph to calculate the probability of an occurrance in that overlap area. No need for erfs.
• i assume that he means that he wants to compute the probability in the green region. That can be done simply by adding the area in the right tail of the first distribution to the ares in the left tail of the second distribution. There is no normalization necessary. – Michael R. Chernick Nov 2 '19 at 23:36
• The separate normalization of the two graphs will change the overlap area unless they have identical SDs. Thus, this answer is generally wrong. – whuber Nov 3 '19 at 13:52 | 2020-04-04T19:36:26 | {
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https://math.stackexchange.com/questions/1620904/how-to-check-convergence-divergence-of-the-series | # How to check convergence/divergence of the series
How to check convergence/divergence of
$\sum \dfrac{\dfrac{1}{2}+(-1)^n}{n}$
I thought in this way that it is of the form $\sum_{i=1}^n \dfrac{1}{2n}+\sum _{i=1}^n \dfrac{(-1)^n}{n}$ i.e. sum of a convergent and a divergent series and hence the original series diverges.
Is it right or there are other alternatives??
I think you have a good idea (with typos) that needs further supporting evidence.
Here's the result you need: Suppose $\sum a_n$ diverges and $\sum b_n$ converges. Then $\sum (a_n + b_n)$ diverges. Proof: Suppose instead that $\sum (a_n + b_n)$ converges. Then $\sum a_n = \sum [(a_n+b_n) - b_n]$ is the difference of two convergent series, hence converges. That's a contradiction, proving the result.
In the problem at hand, we have $a_n = 1/2n, b_n=(-1)^n/n.$ We know $\sum a_n$ diverges, and $\sum b_n$ converges by the alternating series test. Thus your series diverges by the above.
Note that \begin{align} \sum_{n=1}^{\infty}{\frac{\tfrac12+(-1)^n}{n}}&=\sum_{n=1}^{\infty}\left(\frac{\tfrac12+(-1)}{2n-1}+\frac{\tfrac12+(1)}{2n}\right)\\ &=\sum_{n=1}^{\infty}\frac{2n\cdot-\tfrac{1}{2}+(2n-1)\cdot\tfrac32}{2n(2n-1)}\\ &=\sum_{n=1}^{\infty}\frac{2n-\tfrac32}{2n(2n-1)}\\ &\geq\sum_{n=2}^{\infty}\frac{2n-n}{2n(2n-1)}\\ &=\frac12 \sum_{n=2}^{\infty}\frac{1}{2n-1} \end{align} Since the last is smaller than the first sum, and the last sum goes to $+\infty$, so does the first.
Hope this helped!
• How was my solution ;Is it faulty – Learnmore Jan 21 '16 at 11:50
• I think it is correct but you should prove the theorem or reference to a proof if you're using it – vrugtehagel Jan 21 '16 at 11:58
• @Amartya : your argument is correct ; the statement (the sum of a divergent series and a convergent one is divergent is true and easy). – John Steinbeck Jan 21 '16 at 12:13
• @vrugtehagl : I consider the way you write it is not very satisfactory, as you write infinite sums while you didn't prove they make sense. As the term in the sum are not all positive, the meaning of the infinite sum is not always assured. Your proof can certainly be fixed using partial sums and they arguing that your last sum goes to $=\infty$. – John Steinbeck Jan 21 '16 at 12:13 | 2020-07-06T17:56:29 | {
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https://math.stackexchange.com/questions/1919828/prove-that-2n1-gt-n2-for-every-positive-integer-n | # Prove that $2^{n+1} \gt n^2$ for every positive integer $n$.
Prove that $2^{n+1} \gt n^2$ for every positive integer $n$.
I have to prove this using mathematical induction, I know how to start the proof but I get stuck towards the end because I think I either did something wrong or do not know how to continue it.
Here are all of the steps I have done so far:
First I prove that $1 \in S$
This is obviously true because $2^{1+1} \gt 1^2$
Now assume that $a \in S$, we have $2^{a+1} \gt a^2$, now I have to prove that $a+1 \in S$, i.e. $2^{a+2} \gt (a+1)^2$.
$2^{a+2}$ becomes $2 \times 2^{a+1}$ and $(a+1)^2$ becomes $a^2+2a+1$.
We now have $2 \times 2^{a+1} \gt a^2+2a+1$
Multiply both sides of $2^{a+1} \gt a^2$ by 2 and obtain $2 \times 2^{a+1} \gt 2a^2$
to complete the proof all I have to prove is that $2a^2 \gt a^2 +2a +1$
subtracting $a^2$ from both sides I get $a^2 \gt 2a+1$, but this is only true for $a \ge 3$ and I need to prove if for all positive integers.
What did I do wrong, should I have taken a different approach or is there a way to finish what I started?
• Then you can start your induction from $3$. – QZ0 Sep 9 '16 at 0:00
• @ThomasAndrews so I can have two base cases, when $a=1$ and when $a=2$, then prove for $a+1$, can I have two base cases? But how can $a=1$ and $a=2$ work for $2^{a+1} \gt a^2$, but not work for $a^2 \gt 2a+1$ if this is derived from the first inequality? – idknuttin Sep 9 '16 at 0:03
• You can have finite number of base cases. For your second question, you should understand the difference between "sufficient condition" and "necessary condition". – QZ0 Sep 9 '16 at 0:07
You need to prove more base cases.
In the case of $a=1$, we have $2^{1+1} = 4 \gt 1^2 = 1$. In the case of $a=2$, we have $2^{2+1} = 8 \gt 2^2 = 4$.
Assume that for some $a \in \mathbb N, a\ge 3,$ that $2^{a+1} \gt a^2$. Then we must show that the $a+1$ case is true. We have $2^{(a+1)+1)} = 2^{a+1} * 2^1 > 2 * a^2$, since by assumption $2^{a+1} \gt a^2$. Then we have $2^{a+1} * 2^1 > 2 * a^2 > (a+1)^2 = a^2 + 2a + 1$ for $a \ge 3$, and we are done.
• why can we just say that the final part, $a^2 \gt 2a+1$ is obvious and does not require a proof? – idknuttin Sep 9 '16 at 0:16
• Well, you stated in your question that it's only true for $a \ge 3$, so I assumed it was given or perhaps proven earlier. There are many ways of proving this though. Any ideas? – David Bowman Sep 9 '16 at 0:38
• I could probably use mathematical induction again but I don't think that would be necessary, we could move everything to one side and complete the square to get $(a-1)^2 -2 \gt 0$ can we say this is obvious for $a \ge 3$, or is there an easy way of proving it? – idknuttin Sep 9 '16 at 0:46
• You might use the quadratic formula to discover that $f(a) = a^2 - 2a -1$ has zeroes $1- \sqrt{2}$ and $1+ \sqrt{2}$, both of which are less than $3$. Since f is increasing for $a \gt 1, f(a) \gt 0$ for $a \gt 1 +\sqrt{2}$. – David Bowman Sep 9 '16 at 0:52
Another way is using the function $$f(x)=2^{x+1}-x^2$$ whose derivative is $$f'(x)=2^{x+1}\ln 2-2x$$ One has $$2^{x+1}\ln 2-2x\gt 0\iff2^x\ln 2\gt x$$ Hence the derivative is always positive and (besides of increasing) so is the function $f(x)$ for $x\gt 0$ since $f(0)=2$. In particular for $x$ integer positive this is true.
• But then, how do you prove that $2^x\ln(2)>x$? – gniourf_gniourf Dec 4 '16 at 17:16 | 2020-05-25T21:45:38 | {
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https://math.stackexchange.com/questions/1703784/arranging-m-distinct-groups-of-k-elements-each | # Arranging m distinct groups of k elements each?
There are n elements in total which is already divided into m groups of k elements each. Thus, n=m*k.
The question is, how many arrangements of these m groups are possible?
I came to a possible solution of n!/m! but am not sure if that covers all permutations possible among each group (i.e, k! ways of arranging each group).
The original problem is as follows:
There exists an array of n distinct numbers. This array is divided into m sub-sequences of k elements each.
All elements of a subsequence are either greater than or less than all elements of another subsequence.
I have to find the number of possible permutations of the indices in the array which may appear in the sorted array.
• It is unclear whether you are only trying to arrange the $m$ groups or something more complicated. Of course there are $m!$ ways to "arrange" (permute) the $m$ groups. Please make your Question clearer, perhaps by detailing a small example. Mar 19 '16 at 0:17
• What was the original problem you were trying to solve? Mar 19 '16 at 0:17
• Should I change the question to reflect the original problem or add it to the existing problem? Sorry, I am new here. Mar 19 '16 at 0:19
• Please state the original problem you were trying to solve and show the work you did on that problem. In short, we would like to see the context for your current question. Mar 19 '16 at 0:22
• You could also clarify what permutations are allowed. Take a small example for now, $m=n=2$ so there are $4$ elements. Suppose the first group are the capital letters and the second group are the lower case letters. Which of the following are considered the same or distinct or not allowed? ABcd, BAcd, cdAB, cdBA, AcBd,... I.e. must the groups themselves remain together, must the groups be in a specific order, is order within the groups relevant,... Mar 19 '16 at 0:24
The groups are completely determined by the requirement that all elements of a subsequence are either greater than or less than all elements of another subsequence. As you realized, the elements in each of the $m$ groups of $k$ elements can be arranged in $k!$ orders. The $m$ subsequences can be arranged in $m!$ orders. Thus, the number of possible permutations of the indices in the array is $m!(k!)^m$.
• Would this result be strictly greater than $(m.k)!$ or would their relationship depend on values of m and k? Mar 19 '16 at 1:52
• Equality holds when $m = k = 1$. However, in general, $m!(k!)^m \leq (mk)!$ since $(mk)!$ represents all permutations of $mk$ objects, while $m!(k!)^m$ represents all the permutations in which the first $k$ objects are grouped together, the second $k$ objects are grouped together, and so forth. For instance, if $m = 2$ and $k = 3$, the elements in the set $\{1, 2, 3, 4, 5, 6\}$ can be sequenced in $6!$ ways, while the two sets $\{1, 2, 3\}$ and $\{4, 5, 6\}$ can be arranged in $2!$ ways and each can be arranged internally in $3!$ ways for a total of $2!(3!)^2 = 72$ permissible arrangements. Mar 19 '16 at 3:05 | 2021-10-27T17:16:26 | {
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When there are consecutive integers and if their range is equal to the [#permalink]
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When there are consecutive integers and if their range is equal to their median, what is the number of them?
1) The smallest number of them is 5
2) The largest number of them is 15
* A solution will be posted in two days.
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"Only $99 for 3 month Online Course" "Free Resources-30 day online access & Diagnostic Test" "Unlimited Access to over 120 free video lessons - try it yourself" Math Revolution GMAT Instructor Joined: 16 Aug 2015 Posts: 5600 GMAT 1: 800 Q59 V59 GPA: 3.82 Re: When there are consecutive integers and if their range is equal to the [#permalink] ### Show Tags 02 Mar 2016, 22:51 Forget conventional ways of solving math questions. In DS, Variable approach is the easiest and quickest way to find the answer without actually solving the problem. Remember equal number of variables and independent equations ensures a solution. When consecutive integers are there and if their range is equal to their median, what is the number of them? 1) The smallest number of them is 5 2) The largest number of them is 15 In the original condition, there are 2 variables(you need to figure out the starting number and the number of them) and 1 equation(range=median), which should match with the number of equations. So you need 1 variable. For 1) 1 equation, for 2) 1 equation, which is likely to make D the answer. For 1), only 5,6,7,8,9,10,11,12,13,14,15 are possible and so as 2). Therefore, the answer is D. _________________ MathRevolution: Finish GMAT Quant Section with 10 minutes to spare The one-and-only World’s First Variable Approach for DS and IVY Approach for PS with ease, speed and accuracy. "Only$99 for 3 month Online Course"
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Re: When there are consecutive integers and if their range is equal to the [#permalink]
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MathRevolution wrote:
When there are consecutive integers and if their range is equal to their median, what is the number of them?
Let $$a$$ be the first element of the set, $$b$$ be the final element of the set, and $$n$$ be the range.
The median of a set of consecutive numbers is the same as the mean: $$\frac{\text{first term} + \text{last term}}{2}$$
1) The smallest number of them is 5
$$a = 5\\ b = 5 + n\\ \text{median} = \frac{a + b}{2} = \frac{10 + n}{2} = n\\ 2n = 10 + n\\ n = 10\\$$
Sufficient
2) The largest number of them is 15
$$b = 15\\ a = 15 - n\\ \text{median} = \frac{a + b}{2} =\frac{15 + 15 - n}{2} = n\\ 30 - n = 2n\\ 30 = 3n\\ n = 10$$
Sufficient
Knowing $$n$$ and a start point, there are 11 elements from 5 to 15.
(D) each statement alone is sufficient
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Re: When there are consecutive integers and if their range is equal to the [#permalink]
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This seems to be easier if we apply the basic logic for consecutive integers.
We know that for a set of consecutive integers, mean=median.
Consider the first term as 'a' and last term as 'b'.
Range= b-a
Median=Mean= (a+b)/2.
Given: Range=Median or Mean
So, we can say that
b-a = (a+b)/2
b = 3a or a = b/3
Statement 1 : given smallest no is 5 or a = 5
So, b= 3a = 3*5 = 15.
Hence, we get a set {5,6,7,.....,15} which will satisfy all given conditions
Statement 1 is sufficient.
Statement 2: Similarly, given largest no is 15 or b=15
So, a=b/3 = 5
Hence, we get the same set again as {5,.......,14,15}.
Statement 2 is sufficient.
Ans :D
Re: When there are consecutive integers and if their range is equal to the [#permalink] 25 Sep 2017, 21:02
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# When there are consecutive integers and if their range is equal to the
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Powered by phpBB © phpBB Group | Emoji artwork provided by EmojiOne Kindly note that the GMAT® test is a registered trademark of the Graduate Management Admission Council®, and this site has neither been reviewed nor endorsed by GMAC®. | 2018-06-23T02:15:07 | {
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https://www.physicsforums.com/threads/errors-uncertainty-in-measurement-problem.870341/ | # Errors/ Uncertainty in measurement problem
## Homework Statement
[/B]
The surface tension (T) is measured by capillary rise formula $$T = \frac {rh ρg}{2 cos\Theta}$$ . The quantities of ρ, g and θ are taken from the table of constants while the height and diameter are measured as
h = (3.00 + 0.01)cm and
D = (0.250 ± 0.001)cm
Find the percentage error in T
My answer came out to be 1.1 % but doesn't coincide with the answer given (0.7 %).
Thanks for any help!
## Homework Equations
According to the book I refer, the equation for error in product or quotient is -[/B]
$$\frac{\Delta Z} {Z}= \frac {\Delta A} {A} +\frac {\Delta B} {B}$$
where Δz/z, Δa/a and Δb/b are relative errors (letters in deltas present uncertainty in measurement)
## The Attempt at a Solution
I have attempted the solution but have failed to acquire the given answer. Here's what I tried
Acc. to $$\frac{\Delta Z} {Z}= \frac {\Delta A} {A} +\frac {\Delta B} {B}$$
we ignore constants (ρ, Θ and g)
h = (3.00 + 0.01)cm
D = (0.250 ± 0.001)cm
so r = D/2 or r = (0.125 ± 0.001)cm
so
$$\frac{\Delta T} {T}= \frac {\Delta r} {r} +\frac {\Delta h} {h}$$
which gives,
$$\frac{\Delta T} {T} =\frac {0.001} {0.125} +\frac {0.01} {3.00}$$
or,
$$\frac{\Delta T} {T} =0.0080 + 0.0034$$
or,
$$\frac{\Delta T} {T} =0.0113$$
where
$$\frac{\Delta T} {T}$$ is relative error. When we multiply that by 100, we get percent error
therefore percent error = 1.13 %
while the answer given is 0.7 %
Related Introductory Physics Homework Help News on Phys.org
Samy_A
Homework Helper
## Homework Statement
[/B]
The surface tension (T) is measured by capillary rise formula $$T = \frac {rh ρg}{2 cos\Theta}$$ . The quantities of ρ, g and θ are taken from the table of constants while the height and diameter are measured as
h = (3.00 + 0.01)cm and
D = (0.250 ± 0.001)cm
Find the percentage error in T
My answer came out to be 1.1 % but doesn't coincide with the answer given (0.7 %).
Thanks for any help!
## Homework Equations
According to the book I refer, the equation for error in product or quotient is -[/B]
$$\frac{\Delta Z} {Z}= \frac {\Delta A} {A} +\frac {\Delta B} {B}$$
where Δz/z, Δa/a and Δb/b are relative errors (letters in deltas present uncertainty in measurement)
## The Attempt at a Solution
I have attempted the solution but have failed to acquire the given answer. Here's what I tried
Acc. to $$\frac{\Delta Z} {Z}= \frac {\Delta A} {A} +\frac {\Delta B} {B}$$
we ignore constants (ρ, Θ and g)
h = (3.00 + 0.01)cm
D = (0.250 ± 0.001)cm
so r = D/2 or r = (0.125 ± 0.001)cm
so
$$\frac{\Delta T} {T}= \frac {\Delta r} {r} +\frac {\Delta h} {h}$$
which gives,
$$\frac{\Delta T} {T} =\frac {0.001} {0.125} +\frac {0.01} {3.00}$$
or,
$$\frac{\Delta T} {T} =0.0080 + 0.0034$$
or,
$$\frac{\Delta T} {T} =0.0113$$
where
$$\frac{\Delta T} {T}$$ is relative error. When we multiply that by 100, we get percent error
therefore percent error = 1.13 %
while the answer given is 0.7 %
gneill
Mentor
When you divided the diameter by two, the error in the resulting figure should not remain the same as the original. The percentage error of the result should remain the same. So, scale both by the same factor: 0.250/2 = 0.125; 0.001/2 = 0.0005.
When you divided the diameter by two, the error in the resulting figure should not remain the same as the original. The percentage error of the result should remain the same. So, scale both by the same factor: 0.250/2 = 0.125; 0.001/2 = 0.0005.
Thanks, that correction seems to nail it!
One Q - the instrument (say vernier calipers) was used to measure diameter and radius is just half of diameter; then why do we halve the error? Shouldn't the uncertainty remain same ?
gneill
Mentor
Thanks, that correction seems to nail it!
One Q - the instrument (say vernier calipers) was used to measure diameter and radius is just half of diameter; then why do we halve the error? Shouldn't the uncertainty remain same ?
What's the rule when multiplying or dividing by a constant? Does the percent error change?
Dr Transport
Gold Member
You really should be taking the derivative of the original formula with respect to each of the variables and do the root square addition for them to get the uncertainty. Each of the variables in your equation has an uncertainty, even the ones pulled from a reference table and it needs to be included.
Ray Vickson | 2021-03-08T16:09:34 | {
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https://quantumcomputing.stackexchange.com/questions/12604/quantum-katas-tutorials-singlequbitgates-exercise-2-globalphasei | # Quantum Katas - Tutorials - SingleQubitGates - Exercise 2 - GlobalPhaseI
Program.qs
namespace Quantum.Kata.SingleQubitGates {
open Microsoft.Quantum.Intrinsic;
open Microsoft.Quantum.Math;
operation GlobalPhaseI (q : Qubit) : Unit is Adj+Ctl {
X(q);
Z(q);
Y(q);
}
}
Reference.qs
namespace Quantum.Kata.SingleQubitGates {
open Microsoft.Quantum.Intrinsic;
open Microsoft.Quantum.Math;
operation GlobalPhaseI_Reference (q : Qubit) : Unit is Adj+Ctl {
X(q);
Z(q);
Y(q);
}
}
Test.qs
open Microsoft.Quantum.Intrinsic;
open Microsoft.Quantum.Canon;
open Microsoft.Quantum.Diagnostics;
open Microsoft.Quantum.Math;
open Microsoft.Quantum.Convert;
operation ControlledArrayWrapperOperation (op : (Qubit => Unit is Adj+Ctl), qs : Qubit[]) : Unit is Adj+Ctl {
Controlled op([qs[0]], qs[1]);
}
operation AssertEqualOnZeroState (testImpl : (Qubit => Unit is Ctl), refImpl : (Qubit => Unit is Adj+Ctl)) : Unit {
using (qs = Qubit[2]) {
within {
H(qs[0]);
}
apply {
Controlled testImpl([qs[0]], qs[1]);
}
AssertAllZero(qs);
}
}
operation T2_GlobalPhaseI_Test () : Unit {
AssertOperationsEqualReferenced(2, ControlledArrayWrapperOperation(GlobalPhaseI, _), ControlledArrayWrapperOperation(GlobalPhaseI_Reference, _));
}
Driver.cs
using Microsoft.Quantum.Simulation.XUnit;
using Microsoft.Quantum.Simulation.Simulators;
using Xunit.Abstractions;
namespace Quantum.Kata.SingleQubitGates
{
public class TestSuiteRunner
{
public TestSuiteRunner(ITestOutputHelper output)
{
this.output = output;
}
/// <summary>
/// </summary>
[OperationDriver(TestNamespace = "Quantum.Kata.SingleQubitGates")]
public void TestTarget(TestOperation op)
{
using (var sim = new QuantumSimulator())
{
sim.OnLog += (msg) => { output.WriteLine(msg); };
sim.OnLog += (msg) => { Debug.WriteLine(msg); };
op.TestOperationRunner(sim);
}
}
}
}
Yes, this sequence of gates will work to apply the global phase of $$i$$. You can check it using matrix multiplication - a product of matrices $$Y \cdot Z \cdot X$$ will give you a matrix $$\begin{bmatrix} i & 0 \\ 0 & i \end{bmatrix}$$, which corresponds to applying a global phase.
Alternatively, you can implement it using R gate with the phase $$-\pi$$: R(PauliI, -PI(), q);
As a side note, performing measurements this way is not going to detect a global phase introduced by the gate; you need to work with a controlled version of the gate to detect it.
• I tried to work with a controlled version of the gate but I could not add a Entry Potint and run it. How can I run the code with a controlled version ? – theRomanMercury Jun 22 at 19:29
• EntryPoint and controlled version of the gate are completely unrelated. What exactly didn't work for you? You just need to allocate two qubits, prepare the first one in the |+⟩ state and apply controlled gate with the first qubit as control and second as target. – Mariia Mykhailova Jun 23 at 5:38
• I updated my question with your answer. I used this codes from Katas but I did not get an output when I run it. Is it because they are Unit type? – theRomanMercury Jun 23 at 9:03
• Nothing in the code you've pasted writes anything to the output. You can use Message to output classical values and DumpMachine to output internal state of the quantum simulator. – Mariia Mykhailova Jun 23 at 23:10
On IBM Q, you can also use $$U3$$ gate to prepare a global phase operator. $$U3$$ gate is defined as $$U3(\theta, \varphi, \lambda) = \begin{pmatrix} \cos (\theta/2) & -\mathrm{e}^{i\lambda}\sin(\theta/2) \\ \mathrm{e}^{i\varphi}\sin(\theta/2) & \mathrm{e}^{i(\varphi + \lambda)}\cos(\theta/2)\\ \end{pmatrix}.$$ Setting $$\theta = \pi$$ we get $$U3(\pi, \varphi, \lambda) = \begin{pmatrix} 0 & -\mathrm{e}^{i\lambda} \\ \mathrm{e}^{i\varphi} & 0\\ \end{pmatrix}.$$ Let's denote our global phase $$\alpha$$ and set $$\varphi = \alpha$$ and $$\lambda = \alpha + \pi$$. Since $$-\mathrm{e}^{i\pi}=1$$ we have $$U3(\pi, \alpha, \alpha+\pi) = \begin{pmatrix} 0 & \mathrm{e}^{i\alpha} \\ \mathrm{e}^{i\alpha} & 0\\ \end{pmatrix},$$ which is $$\mathrm{e^{i\alpha}}X$$. To get $$\mathrm{e^{i\alpha}}I$$, we apply another $$X$$ gate.
So, global phase gate is implemented as $$X\,\,U3(\pi,\alpha,\alpha+\pi)$$, where $$\alpha$$ is global phase.
EDIT (solution in Q#)
A $$R1$$ gate is defined as $$R1(\theta) = \begin{pmatrix} 1 & 0 \\ 0 & \mathrm{e}^{i\theta} \end{pmatrix}$$
An operation $$X\,R1(\theta)$$ is described by matrix $$\begin{pmatrix} 0 & \mathrm{e}^{i\theta} \\ 1 & 0 \end{pmatrix}$$
If we apply this operation twice (i.e. $$[X\,R1(\theta)]^2$$ ), we get
$$\begin{pmatrix} \mathrm{e}^{i\theta} & 0 \\ 0 & \mathrm{e}^{i\theta} \end{pmatrix},$$
which is a global phase gate with arbitrary phase $$\theta$$.
So, global phase gate in Q# can be realized as $$[X\,R1(\theta)]^2$$.
• I am focusing on Q# language, thanx a lot for your answer. – theRomanMercury Jun 23 at 9:04
• @theRomanMercury: Sorry, I did not take this into account. Please find solution (global phase gate with arbitrary phase) with $R1$ and $X$ gates for Q#. – Martin Vesely Jun 24 at 7:34
I apologize, but I can not comment on @Martin Vesely answer above. I am glad that my amendments have already been adopted, I’ll leave only a short note, in case someone is interested in this: In Qiskit programs (e.g. see) I use X before U3 in my variants of phase shift circuits, in an attempt to avoid "distortion" during further processing of Qiskit (for example, by transpiler):
qc.x(qubit)
qc.u3(np.pi, gamma, np.pi + gamma, qubit)
If not for similar considerations, then the order here would not be important (as in other similar cases, e.g. $$[X\,R1(\theta)]^2 = [R1(\theta)\,X]^2$$) | 2020-10-25T22:21:32 | {
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https://waqsp.org/telescope-price-zwgol/73f536-extreme-value-theorem-open-interval | This example was to show you the extreme value theorem. If has an extremum on an open interval , then the extremum occurs at a critical point. History. Critical points are determined by using the derivative, which is found with the Chain Rule. If f(x) exists for all values of x in the open interval (a,b) and f has a relative extremum at c, where a < c < b, then if f0(c) exists, then f0(c) = 0. Proof of the Extreme Value Theorem Theorem: If f is a continuous function defined on a closed interval [a;b], then the function attains its maximum value at some point c contained in the interval. The Extreme Value Theorem (EVT) says: If a function f is continuous on the closed interval a ≤ x ≤ b , then f has a global minimum and a global maximum on that interval. point. know that the absolute extremes occur at either the endpoints, x=0 and x = 3, or the The function Proof of the Extreme Value Theorem Theorem: If f is a continuous function defined on a closed interval [a;b], then the function attains its maximum value at some point c contained in the interval. Open interval. We use the logarithm to compute the derivative of a function. Closed interval domain, … Bolzano's proof consisted of showing that a continuous function on a closed interval was bounded, and then showing that the function attained a maximum and a minimum value. So is a value attained at least three times: inside the open interval (,) , inside of (, ~) and inside of (~, ~). If the interval is open or the function has even one point of discontinuity, the function may not have an absolute maximum or absolute minimum over For example, consider the functions shown in (Figure) (d), (e), and (f). Correspondingly, a metric space has the Heine–Borel property if every closed and bounded set is also compact. Are you sure you want to do this? It states that if f is a continuous function on a closed interval [a, b], then the function f has both a minimum and a maximum on the interval. In this section we learn the definition of continuity and we study the types of The absolute maximum is shown in red and the absolute minimumis in blue. • Extreme Value Theorem: • A function can have only one absolute maximum or minimum value (y-coordinate), but it can occur at more than one x-coordinate. The Inverse Function Theorem (continuous version) 11. numbers of f(x) in the interval (0, 3). Since the endpoints are not included, they can't be the global extrema, and this interval has no global minimum or maximum. The absolute extremes occur at either the endpoints, x=\text {-}1, 6 or the critical In order to use the Extreme Value Theorem we must have an interval that includes its endpoints, often called a closed interval, and the function must be continuous on that interval. This is used to show thing like: There is a way to set the price of an item so as to maximize profits. Inequalities and behaviour of f(x) as x →±∞ 17. Solution: The function is a polynomial, so it is continuous, and the interval is closed, Play this game to review undefined. This video explains the Extreme Value Theorem and then works through an example of finding the Absolute Extreme on a Closed Interval. The Extreme Value Theorem 10. at a Regular Point of a Surface. function f(x) yields: The absolute maximum is \answer {2e^4} and it occurs at x = \answer {2}.The absolute minimum is \answer {-1/(2e)} and it occurs at x = \answer {-1/2}. Solution: First, we find the critical numbers of f(x) in the interval [\text {-}1, 6]. In this section we discover the relationship between the rates of change of two or Since is compact, The Mean Value Theorem states that the rate of change at some point in a domain is equal to the average rate of change of that domain. An important Theorem is theExtreme Value Theorem. fails to hold, then f(x) might fail to have either an absolute max or an absolute min If f'(c) is undefined then, x=c is a critical number for f(x). However, for a function defined on an open or half-open interval… Extreme Value Theorem Theorem 1 below is called the Extreme Value theorem. Extreme Value Theorem: When examining a function, we may find that it it doesn't have any absolute extrema. It ... (-2, 2), an open interval, so there are no endpoints. Lagrange mean value theorem; Bound relating number of zeros of function and number of zeros of its derivative; Facts used. y = x2 0 ≤ x ≤2 y = x2 0 ≤ x ≺2 4.1 Extreme Values of Functions Day 2 Ex 1) A local maximum value occurs if and only if f(x) ≤ f(c) for all x in an interval. State whether the function has absolute extrema on its domain $$g(x)= \begin{cases} x & x> 0 \\ 3 & x\leq 0 \end{cases}$$ Join the initiative for modernizing math education. Thus, a bound of infinity must be an open bound. Plugging these special values into the original function f(x) yields: The absolute maximum is \answer {17} and it occurs at x = \answer {-2}.The absolute minimum is \answer {-15} and it occurs at x = \answer {2}. This theorem is sometimes also called the Weierstrass extreme value theorem. 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The Extreme Value Theorem guarantees both a maximum and minimum value for a function under certain conditions. This becomes 3x^2 -12x= 0 First, since we have a closed interval (i.e. I know it must be continuous for the interval, but must it be closed? Suppose that f(x) is defined on the open interval (a,b) and that f(x) has an absolute max at x=c. Extreme Value Theorem: f is continuous on [a;b], then f has an absolute max at c and and absolute min at d, where c;d 2[a;b] 7. Intermediate Value Theorem and we investigate some applications. Portions of this entry contributed by John We learn to compute the derivative of an implicit function. If has an absolute maximum or absolute minimum at a point in the interval, then is a critical number for . Determine all critical points in the given interval and evaluate the function at these critical points and at the endpoints of the interval 3. Establish that the function is continuous on the closed interval 2. Closed Interval Method • Test for absolute extrema • Only use if f (x) is continuous on a closed interval [a, b]. (If the max/min occurs in more than one place, list them in ascending order).The absolute maximum is \answer {4} and it occurs at x = \answer {-2} and x=\answer {1}. The absolute maximum is \answer {0} and it occurs at x = \answer {-2}. Below, we see a geometric interpretation of this theorem. I know it's pretty vital for the theorem to be able to show the values of f(a) and f(b), but what if I have calculated the limits as x -> a (from the right hand side) and x -> b (from the left hand side)? Use the differentiation rules to compute derivatives. In this section we will solve the problem of finding the maximum and minimum We learn how to find the derivative of a power function. When you do the problems, be sure to be aware of the difference between the two types of extrema! Proof: There will be two parts to this proof. We solve the equation f'(x) =0. interval. In this section we interpret the derivative as an instantaneous rate of change. To find the relative extrema of a function, you first need to calculate the critical values of a function. This has two important corollaries: . Wolfram Web Resource. Walk through homework problems step-by-step from beginning to end. For example, let’s say you had a number x, which lies somewhere between zero and 100: The open interval would be (0, 100). 3. The absolute minimum is \answer {0} and it occurs at x = \answer {-3} and x=\answer {0}. compute the derivative of an area function. We determine differentiability at a point. Notice how the minimum value came at "the bottom of a hill," and the maximum value came at an endpoint. For the extreme value theorem to apply, the function must be continuous over a closed, bounded interval. A description of how to find the arguments that maximize and minimize the value of a function that is continuous on a closed interval. © 2013–2021, The Ohio State University — Ximera team, 100 Math Tower, 231 West 18th Avenue, Columbus OH, 43210–1174. is a polynomial, so it is differentiable everywhere. Unlimited random practice problems and answers with built-in Step-by-step solutions. The next theorem is called Rolle’s Theorem and it guarantees the existence of an extreme value on the interior of a closed interval, under certain conditions. This theorem is called the Extreme Value Theorem. 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If a function is continuous on a closed The first derivative can be used to find the relative minimum and relative maximum values of a function over an open interval. In this section we use definite integrals to study rectilinear motion and compute itself be compact. In order to use the Extreme Value Theorem we must have an interval that includes its endpoints, often called a closed interval, and the function must be continuous on that interval. We don’t want to be trying to find something that may not exist. The closed interval—which includes the endpoints— would be [0, 100]. A local minimum value … However, if that interval was an open interval of all real numbers, (0,0) would have been a local minimum. Also note that while $$0$$ is not an extreme value, it would be if we narrowed our interval to $$[-1,4]$$. Using the Extreme Value Theorem 1. (a,b) as opposed to [a,b] Also note that while $$0$$ is not an extreme value, it would be if we narrowed our interval to $$[-1,4]$$. 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THE EXTREME-VALUE THEOREM (EVT) 27 Interlude: open and closed sets We went about studying closed bounded intervals… Here is a detailed, lecture style video on the Extreme Value Theorem: Calculus I, by Andrew In this section we learn the definition of the derivative and we use it to solve the Extreme Value Theorem If f is continuous on a closed interval [a,b], then f has both a maximum and minimum value. The derivative is f'(x) = 2x - 4 which exists for all values of x. is differentiable on the open interval , i.e., the derivative of exists at all points in the open interval .. Then, there exists in the open interval such that . First, we find the critical Hence f'(c) = 0 and the theorem is In this section we learn to compute the value of a definite integral using the A local minimum value … In this section we compute derivatives involving. The extreme value theorem gives the existence of the extrema of a continuous function defined on a closed and bounded interval. knowledge of derivatives. In this section, we use the derivative to determine intervals on which a given function Depending on the setting, it might be needed to decide the existence of, and if they exist then compute, the largest and smallest (extreme) values of a given function. In mathematical analysis, the intermediate value theorem states that if f is a continuous function whose domain contains the interval [a, b], then it takes on any given value between f(a) and f(b) at some point within the interval.. This video explains the Extreme Value Theorem and then works through an example of finding the Absolute Extreme on a Closed Interval. The extreme value theorem gives the existence of the extrema of a continuous function defined on a closed and bounded interval. In this section we examine several properties of the indefinite integral. We solve the equation f'(x) =0. f'(x) =0. If the interval $$I$$ is open or the function has even one point of discontinuity, the function may not have an absolute maximum or absolute minimum over $$I$$. Try the following: The first graph shows a piece of a parabola on a closed interval. Below, we see a geometric interpretation of this theorem. Basically Rolle ‘s theorem states that if a function is differentiable on an open interval, continuous at the endpoints, and if the function values are equal at the endpoints, then it has at least one horizontal tangent. Extreme Value Theorem: If a function is continuous in a closed interval , with the maximum of at and the minimum of at then and are critical values of Proof: The proof follows from Fermat’s theorem and is left as an exercise for the student. y = x2 0 ≤ x ≤2 y = x2 0 ≤ x ≺2 4.1 Extreme Values of Functions Day 2 Ex 1) A local maximum value occurs if and only if f(x) ≤ f(c) for all x in an interval. average value. This is a good thing of course. The procedure for applying the Extreme Value Theorem is to first establish that the function is continuous on the closed interval. Plugging these special values into the original However, if that interval was an open interval of all real numbers, (0,0) would have been a local minimum. View Chapter 3 - Limits.pdf from MATHS MA131 at University of Warwick. So, it is (−∞, +∞), it cannot be [−∞, +∞]. If has an extremum the point of tangency. analysis includes the position, velocity and acceleration of the particle. It is not de ned on a closed interval, so the Extreme Value Theorem does not apply. Thus f'(c) \geq 0. function. In this section we compute limits using L’Hopital’s Rule which requires our For example, [0,1] means greater than or equal to 0 and less than or equal to 1. In calculus, the extreme value theorem states that if a real-valued function f is continuous on the closed interval [a,b], then f must attain a maximum and a minimum, each at least once.That is, there exist numbers c and d in [a,b] such that: three step process. A set $${\displaystyle K}$$ is said to be compact if it has the following property: from every collection of open sets $${\displaystyle U_{\alpha }}$$ such that $${\textstyle \bigcup U_{\alpha }\supset K}$$, a finite subcollection $${\displaystyle U_{\alpha _{1}},\ldots ,U_{\alpha _{n}}}$$can be chosen such that $${\textstyle \bigcup _{i=1}^{n}U_{\alpha _{i}}\supset K}$$. • Three steps polynomial, so it is differentiable everywhere. This becomes x^3 -6x^2 + 8x = 0 and the solutions are x=0, x=2 and x=4 (verify). The Extreme Value Theorem. In this section we learn to reverse the chain rule by making a substitution. The idea that the point $$(0,0)$$ is the location of an extreme value for some interval is important, leading us to a definition. Fermat’s Theorem Suppose is defined on the open interval . Example . The quintessential point is this: on a closed interval, the function will have both minima and maxima. Extreme Value Theorem If a function f {\displaystyle f} is continuous on a closed interval [ a , b ] {\displaystyle [a,b]} then there exists both a maximum and minimum on the interval. A continuous function ƒ (x) on the closed interval [a,b] showing the absolute max (red) and the absolute min (blue).. It is also important to note that the theorem tells us that the The extreme value theorem states that if is a continuous real-valued function on the real-number interval defined by , then has maximum and minimum values on that interval, which are attained at specific points in the interval. That makes sense. In such a case, Theorem 1 guarantees that there will be both an absolute maximum and an absolute minimum. max and the min occur in the interval, but it does not tell us how to find If a continuous function has values of opposite sign inside an interval, then it has a root in that interval (Bolzano's theorem). Solution: First, we find the critical numbers of f(x) in the interval [\text {-}1, 3]. The absolute minimum is \answer {-27} and it occurs at x = \answer {1}. Practice online or make a printable study sheet. The image below shows a continuous function f(x) on a closed interval from a to b. Hints help you try the next step on your own. This is what is known as an existence theorem. 2. •Note: If the interval is open, then the endpoints are. A point is considered a minimum point if the value of the function at that point is less than the function values for all x-values in the interval. endpoints, x=-1, 2 or the critical number x = -1/3. In this section we use properties of definite integrals to compute and interpret If f'(c) is defined, then The first is that f(x) is Play this game to review undefined. That is, there exist x 1,x 2 ∈ [a,b] so that f(x 1) = m and f(x 2) = M, and m ≤ f(x) ≤ M for all x ∈ [a,b]. Extreme Value Theorem. The largest and smallest values from step two will be the maximum and minimum values, respectively from the definition of the derivative we have f'(c) = \lim _{h \to 0^-} \frac {f(c+h) -f(c)}{h} = \lim _{h \to 0^-} \frac {f(c+h) -f(c)}{h} The difference quotient in the left If a function is continuous on a closed interval , then has both a maximum and a minimum on . Then, for all >0, there is an algorithm that computes a price vector whose revenue is at least a (1 )-fraction of the optimal revenue, and whose running time is polynomial in max ˆ the equation f'(x) =0 gives x=2 as the only critical number of the function. When moving from the real line $${\displaystyle \mathbb {R} }$$ to metric spaces and general topological spaces, the appropriate generalization of a closed bounded interval is a compact set. The Extreme Value Theorem ... as x !1+ there is an open circle, so the lower bound of y = 1 is approached but not attained . Preview this quiz on Quizizz. To find the relative extrema of a function, you first need to calculate the critical values of a function. This is usually stated in short as "every open cover of $${\displaystyle K}$$ has a finite subcover". which has two solutions x=0 and x = 4 (verify). Solution: First, we find the critical numbers of f(x) in the interval [-1, 0]. critical number x = 2. (or both). Extreme Value Theorem. discontinuities. This example was to show you the extreme value theorem. Remark: An absolute extremum of a function continuous on a closed interval must either be a relative extremum or a function value at an endpoint of the interval. The absolute extremes occur at either the interval around c (open interval around c means that the immediate values to the left and to the right of c are in that open interval) 6. How would you like to proceed? right hand limit is negative (or zero) and so f'(c) \leq 0. https://mathworld.wolfram.com/ExtremeValueTheorem.html, Normal Curvature In this section we learn to compute general anti-derivatives, also known as indefinite If the interval is open or the function has even one point of discontinuity, the function may not have an absolute maximum or absolute minimum over For example, consider the functions shown in … Proof: There will be two parts to this proof. The function is a Noting that x = 4 is not in Regardless, your record of completion will remain. In this lesson we will use the tangent line to approximate the value of a function near integrals. Extreme Value Theorem If is continuous on the closed interval , then there are points and in , such that is a global maximum and is a global minimum on . values of a continuous function on a closed interval. and x = 2. In this section we learn about the two types of curvature and determine the curvature In order to utilize the Mean Value Theorem in examples, we need first to understand another called Rolle’s Theorem. The main idea is finding the location of the absolute max and absolute min of a 4 Extreme Value Theorem If f is continuous on a closed interval a b then f from MATH 150 at Simon Fraser University tangent line problem. Vanishing Derivative Theorem Assume f(x) is a continuous function defined on an open interval (a,b). be compact. interval , so it must Extreme value theorem We compute the derivative of a composition. Extreme Value Theorem If f is continuous on a closed interval [a,b], then f has both a maximum and minimum value. But the difference quotient in the It is important to note that the theorem contains two hypothesis. If the function f is continuous on the closed interval [a,b], then f has an absolute maximum value and an absolute minimum value on [a,b]. You cannot have a closed bound of ±∞ because ∞ is never a value that can actually be reached. There are a couple of key points to note about the statement of this theorem. From MathWorld--A The extreme value theorem was originally proven by Bernard Bolzano in the 1830s in a work Function Theory but the work remained unpublished until 1930. • Three steps/labels:. The absolute extremes occur at either the endpoints, x=\text {-}1, 3 or the critical The #1 tool for creating Demonstrations and anything technical. The Extreme Value Theorem ... as x !1+ there is an open circle, so the lower bound of y = 1 is approached but not attained . Terminology. maximum and a minimum on proved. An open interval does not include its endpoints, and is indicated with parentheses. This theorem is sometimes also called the Weierstrass extreme value theorem. Fermat’s Theorem. . The Weierstrass Extreme Value Theorem. Using the product rule and the chain rule, we have f'(x) = 1\cdot e^{3x} + x\cdot e^{3x} \cdot 3 which simplifies to (3x+1)e^{3x}. Hence f(x) has two critical numbers in the Which of the following functions of x is guaranteed by the Extreme Value Theorem to have an absolute maximum on the interval [0, 4]? Our so by the Extreme Value Theorem, we know that this function has an absolute No maximum or minimum values are possible on the closed interval, as the function both increases and decreases without bound at x = π/2. However, it never reaches the value of $0.$ Notice that this function is not continuous on a closed bounded interval containing 0 and so the Extreme Value Theorem does not apply. Explore anything with the first computational knowledge engine. In this section we prepare for the final exam. Continuous, 3. For the extreme value theorem to apply, the function must be continuous over a closed, bounded interval. The Heine–Borel theorem asserts that a subset of the real line is compact if and only if it is both closed and bounded. and the denominator is negative. within a closed interval. Derivatives of sums, products and composites 13. function. interval around c (open interval around c means that the immediate values to the left and to the right of c are in that open interval) 6. In this section we learn a theoretically important existence theorem called the them. Fermat’s Theorem Suppose is defined on the open interval . These values are often called extreme values or extrema (plural form). Thus f(c) \geq f(x) for all x in (a,b). Real-valued, 2. Differentiation 12. of functions whose derivatives are already known. The Extreme Value Theorem says: If a function f is continuous on the closed interval a ≤ x ≤ b, then f has a global minimum and a global maximum on that interval. https://mathworld.wolfram.com/ExtremeValueTheorem.html. The standard proof of the first proceeds by noting that is the continuous image of a compact set on the We learn the derivatives of many familiar functions. more related quantities. Since we know the function f(x) = x2 is continuous and real valued on the closed interval [0,1] we know that it will attain both a maximum and a minimum on this interval. The derivative is 0 at x = 0 and it is undefined at x = -2 and x = 2. We will also determine the local extremes of the at the critical number x = 0. Local Extrema, critical points, Rolle’s Theorem 15. The Mean Value Theorem 16. If we don’t have a closed interval and/or the function isn’t continuous on the interval then the function may or may not have absolute extrema. The Extreme value theorem requires a closed interval. If you update to the most recent version of this activity, then your current progress on this activity will be erased. Theorem 2. You are about to erase your work on this activity. Then f has both a Maximum and a Minimum value on [a,b].#Extreme value theorem (The circle, in fact.) So the extreme value theorem tells us, look, we've got some closed interval - I'm going to speak in generalities here - so let's say that's our X axis and let's say we have some function that's defined on a closed interval. Hence f(x) has one critical numbers x = 0,2. The quintessential point is this: on a closed interval, the function will have both minima and maxima. On a closed interval, always remember to evaluate endpoints to obtain global extrema. this critical number is in the interval (0,3). In this section we analyze the motion of a particle moving in a straight line. Extreme Value Theorem If a function f {\displaystyle f} is continuous on a closed interval [ a , b ] {\displaystyle [a,b]} then there exists both a maximum and minimum on the interval. In order for the extreme value theorem to be able to work, you do need to make sure that a function satisfies the requirements: 1. fundamental theorem of calculus. For example, (0,1) means greater than 0 and less than 1.This means (0,1) = {x | 0 < x < 1}.. A closed interval is an interval which includes all its limit points, and is denoted with square brackets. the interval [\text {-}1,3] we see that f(x) has two critical numbers in the interval, namely x = 0 If we don’t have a closed interval and/or the function isn’t continuous on the interval then the function may or may not have absolute extrema. 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extreme value theorem open interval 2021 | 2021-07-27T10:49:08 | {
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https://math.stackexchange.com/questions/1788692/topological-finer-and-separability | # Topological finer and separability
I have this question:
Let $(X, d_1)$ and $(X,d_2)$ be two metric spaces. Suppose $d_1$ is topologically finer than $d_2$. What is the relationship between these two statements?
(i) $(X,d_1)$ is separable.
(ii) $(X,d_2)$ is separable.
I think the implication (i) $\Rightarrow$ (ii) holds, but not the other way around. I was trying to prove it:
Suppose $(X,d_1)$ is separable. Then $X$ has an at most countable dense subset for the $d_1$ metric. So there exists an at most countable subset $A \subset X$ such that $\overline{A} = X$. Since $d_1$ is topologically finer than $d_2$, every open subset for the $d_2$ metric is also open for the $d_1$ metric. This is equivalent in saying that every closed $d_2$-subset is also $d_1$-closed.
This is as far I could go. What can I conclude from this? And am I also correct in saying the other implication is false?
• Sorry, I misread. You are correct is thinking that separability in a finer topology implies separability in a coarser topology. I have deleted temporarily my answer and will repair. That's done now. I hope I did not cause any confusion. – drhab May 17 '16 at 9:30
• I saw your answer, and I was thinking you confused $d_1$ with $d_2$. But your answer still helped me a lot in giving me the right idea! – Kamil May 17 '16 at 9:36
• Nice to hear that. There was indeed some confusion by me about $d_1,d_2$. – drhab May 17 '16 at 9:38
Counterexample of (ii)$\implies$(i):
Let $d_1$ be defined on $\mathbb R$ by $d_1(x,y)=0$ if $x=y$ and $d(x,y)=1$ otherwise.
Let $d_2$ be defined on $\mathbb R$ by $d_2(x,y)=|x-y|$.
Then $\mathbb R$ equipped with $d_1$ results in the discrete topology (finest as possible), hence no separability on the uncountable $\mathbb R$.
On the other hand $\mathbb R$ equipped with $d_2$ results in the usual topology, hence separability.
Proof of (i)$\implies$(ii):
Let $X$ equipped with some topology $\tau$ be separable.
Then a countable and dense $S\subseteq X$ exists.
Density of $S$ means exactly that $U\cap S\neq\varnothing$ for each $U\in\tau$.
If $\sigma$ is another topology on $X$ and $\tau$ is finer then also $U\cap S\neq\varnothing$ for $U\in\sigma$. This because $\sigma\subseteq\tau$.
So again $S$ is dense. Of course it has not lost its countability and we conclude that $X$ equipped with the coarser $\sigma$ is separable.
A counterexample for the opposite: We use the set $\Bbb R$ with usual metric $d_2(x,y)=|x-y|$ and define $$d_{1}\left(x,y\right)=\cases{ 0&\text{if }x=y\cr 1&\text{if }x\neq y}$$ We know that $d_2$ make a separable space (use $\Bbb Q$). $d_1$, on the other hand, yields discrete topology, so the only dense set is the whole space, making it not separable.
$d_1$'s topology is finer than $d_2$'s as it's discrete and therefore finer than any other topology.
• $x \not = y$ \not = – Olba12 May 17 '16 at 8:50
• Ah, I see thanks. But is $d_1$ topologically finer than $d_2$ ? – Kamil May 17 '16 at 8:54
• @Kamil Every set open in $d_2$ (in fact every set) is open in $d_1$. – Heimdall May 17 '16 at 8:59
• @Olba12 When I was typing the answer, I tried different combinations: \ne, \neq, \not=, \not =; the preview didn't render either of them correctly. This question uses \neq, I'm quoting the source: If det $A\neq0$, the matrix, so what's wrong with \neq and other possible no-equals in my answer? Maybe it works in question and comments but not in answers? I'm asking on meta. – Heimdall May 17 '16 at 9:46
• Weird, \not = has always been working for me... haha – Olba12 May 17 '16 at 10:08 | 2019-07-16T14:24:42 | {
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https://mathhelpboards.com/threads/number-of-ways-to-divide-a-square-grid-in-half.26269/ | Number of ways to divide a square grid in half
HapaxOromenon
New member
Consider a dotted square grid with $4$ rows and $4$ columns, as shown:
By drawing a series of connected straight lines from dot to dot, it is possible to divide the square (which effectively has side length $3$) into two parts of equal area. One way of doing this is as shown:
Other examples of ways of doing this are: ; ; and .
The question I am investigating is how many such ways there are in total, or if such a calculation is too difficult, I would at least like to have a systematic way of listing all of the ways. But I'm not sure really where to start. Any assistance would be greatly appreciated.
Opalg
MHB Oldtimer
Staff member
Consider a dotted square grid with $4$ rows and $4$ columns, as shown:
By drawing a series of connected straight lines from dot to dot, it is possible to divide the square (which effectively has side length $3$) into two parts of equal area. One way of doing this is as shown:
Other examples of ways of doing this are: ; ; and .
The question I am investigating is how many such ways there are in total, or if such a calculation is too difficult, I would at least like to have a systematic way of listing all of the ways. But I'm not sure really where to start. Any assistance would be greatly appreciated.
Hi HapaxOromenon , and welcome to MHB!
This looks like an interesting problem, and I'm not sure whether it has a straightforward solution. But first, we need to be sure what you are trying to count. You want to partition the square grid into two parts, and you say that these parts should have the same area. In all the examples that you show, the two parts are congruent. So my first question is this: Would you allow a partition like this
\begin{tikzpicture}
[scale=2]\foreach \i in {0,...,3}
{
\draw[thin,dotted] (\i,0) -- (\i,3) ;
\draw[thin,dotted] (0,\i) -- (3,\i) ;
\foreach \j in {0,...,3}
\fill [black,opacity=.5] (\i,\j) circle (2pt) ; }
\draw[thick] (0,2) -- (2,2) -- (2,1) -- (3,0) ;
\end{tikzpicture}
where the two parts have the same area (4.5 square units) but they have different shapes?
The second question is whether you would consider two partitions to be the same if one of them is just a rotation or reflection of the other one. For example, do you consider these partitions
\begin{tikzpicture}
[scale=2]\foreach \i in {0,...,3}
{
\draw[thin,dotted] (\i,0) -- (\i,3) ;
\draw[thin,dotted] (0,\i) -- (3,\i) ;
\foreach \j in {0,...,3}
\fill [black,opacity=.5] (\i,\j) circle (2pt) ; }
\draw[thick] (0,2) -- (1,2) -- (2,1) -- (3,1) ;
\end{tikzpicture} ................ \begin{tikzpicture}
[scale=2]\foreach \i in {0,...,3}
{
\draw[thin,dotted] (\i,0) -- (\i,3) ;
\draw[thin,dotted] (0,\i) -- (3,\i) ;
\foreach \j in {0,...,3}
\fill [black,opacity=.5] (\i,\j) circle (2pt) ; }
\draw[thick] (2,3) -- (2,2) -- (1,1) -- (1,0) ;
\end{tikzpicture}
to be the same, or different?
HapaxOromenon
New member
1) Yes, I would count this as a valid cut.
2) I would consider them to be different.
Opalg
MHB Oldtimer
Staff member
Each cut must start and finish at points on the boundary of the grid. There are two sorts of boundary points, namely the four corners of the grid, and the eight points on the edges that are not corners (I'll call these edge points).
Trying to find a systematic way to count the number of possible cuts, I split them into three classes:
1. cuts that go from corner to corner;
2. cuts that go from a corner to an edge point;
3. cuts that go between two edge points.
For the first class, there are nine cuts between the northwest and southeast corners of the grid, as in this diagram
.
Rotating those grids through a right angle, you get the same number of cuts between the northeast and southwest corners of the grid. So altogether there are 18 cuts of class 1.
For the second class, there are ten cuts starting at the northwest corner of the grid and ending at an edge point, as in this diagram
.
Rotating them, you get the same number of cuts starting at each of the other corners of the grid. So there are 40 cuts of class 2.
That leaves class 3, the cuts going from an edge point to a point on the opposite edge. I found 17 cuts starting from the point immediately below the northwest corner of the grid, as in this diagram
Flipping each of these about a horizontal axis, you get the same number of cuts starting from the point immediately above the southwest corner of the grid. So there are 34 cuts of class 3 going from side to side of the grid. Rotating them through a right angle, you get another 34 cuts of class 3 between the top and bottom of the grid. So altogether there are 68 cuts of class 3.
Thus the overall total number of cuts is 18 + 40 + 68 = 126 (unless I have overlooked any other possible cuts). | 2020-07-06T06:01:50 | {
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https://math.stackexchange.com/questions/2745825/determine-the-constant-alpha-in-mathbbr-and-the-vector-w-in-mathbbr/2746416 | Determine the constant $\alpha \in \mathbb{R}$ and the vector $w \in \mathbb{R}^3$ that is perpendicular to the vector u such that v = α u + w.
Determine the constant $\alpha \in \mathbb{R}$ and the vector $w \in \mathbb{R}^3$ that is perpendicular to the vector $u$ such that $v = \alpha u + w$.
Need a little help with this question from my linear algebra course, not sure how to find the solution. I know that that the dot product of w and u is equal to zero.
the vectors give are $u = (2, -2, 1) v = (3, 1, -1)$
• So use the dot product on both sides of the $v$ equation... – N74 Apr 20 '18 at 10:21
• Sorry I am still a little confused. use the dot with what on the left side of the equation? – Darragh O'Flaherty Apr 20 '18 at 10:27
• On the left side you will have $v \cdot u$ – N74 Apr 20 '18 at 10:30
Let $w=(w_1,w_2,w_3)$. If $w$ is perpendicular to $u$, the equation $w\cdot u=0$ must be satisfied, that is, $2w_1-2w_2+w_3=0$. As you see there are lots of solutions.
If you want $v=\alpha u+w$, then you have another equation, in this cases a system of equations indeed, namely
$3=2\alpha +w_1$
$1=2\alpha+w_2$
$-1=\alpha+w_3$
Adding the equation $2w_1-2w_2+w_3=0$ from above you've got a linear system of 4 equations and 4 unknowns, so there is precisely one solution that I'm sure you're able to find.
• Perfect thank you very much, just the guidance I needed – Darragh O'Flaherty Apr 20 '18 at 11:28
You have $v= \alpha u + w$. Now, dot product by $u$ gives $$v \cdot u = \alpha u \cdot u + w \cdot u$$ We know $w \cdot u=0$ so the equation gives $$3 =\alpha 9$$ so $$\alpha={ 1 \over 3}$$. Knowing $\alpha$ it is straightforward that $$w=\left ({7 \over 3}, {5 \over 3}, -{4 \over 3} \right )$$ | 2021-04-22T17:06:35 | {
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http://mathoverflow.net/questions/20275/maps-inducing-zero-on-homotopy-groups-but-are-not-null-homotopic | # Maps inducing zero on homotopy groups but are not null-homotopic
Today my fellow grad student asked me a question, given a map f from X to Y, assume $f_*(\pi_i(X))=0$ in Y, when is f null-homotopic?
I search the literature a little bit, D.W.Kahn
http://projecteuclid.org/DPubS/Repository/1.0/Disseminate?view=body&id=pdf_1&handle=euclid.pjm/1102995805
And M.Sternstein has worked on this, and Sternstein even got a necessary and sufficient condition, for suitable spaces.
http://www.jstor.org/stable/pdfplus/2037939.pdf
However, his condition is a little complicated for me as a beginner. Right now I just wanted a counter example of a such a map. Kahn in his paper said one can have many such examples using Eilenberg Maclance spaces. Well, we can certainly show a lot of map between E-M spaces induce zero map on homopoty groups just by pure group theoretic reasons, but I can not think of a easy example when you can show that map, if it exists, is not null-homotopic. Could someone give me some hint?
or, maybe even some examples arising from manifolds?
-
Consider ordinary singular cohomology with varying coefficients. You can look at the short exact sequence of abelian groups:
$$0 \to \mathbb{Z}/2 \to \mathbb{Z}/4 \to \mathbb{Z}/2 \to 0$$
This gives rise, for any space X, to a short exact sequence of chain complexes:
$$0 \to C^i(X;\mathbb{Z}/2) \to C^i(X;\mathbb{Z}/4) \to C^i(X;\mathbb{Z}/2) \to 0$$
and hence you get a long exact sequence in cohomology. Thus we get an interesting boundary map known as the Bockstein
$$H^i(X; \mathbb{Z}/2) \to H^{i+1}(X; \mathbb{Z}/2).$$
This is natural in X and so is represented by a (homotopy class of) map(s) of Eilenberg-Maclane spaces:
$$K(i, \mathbb{Z}/2) \to K(i+1, \mathbb{Z}/2)$$
This map is necessarily zero on homotopy groups. To show that this map is not null-homotopy, you just need to find a space for which the Bockstein is non-trivial. There are lots of examples of this. Rather then explain one, I suggest you look up "Bockstein homomorphism" in a standard algebraic topology reference, e.g. Hatcher's book.
-
Thanks, this argument is very nice! I guess it is not very hard to find a space with non trivial Bockstein, e.g.$M=RP^2$, we have $H^i(M,Z/2)=H^i(M,Z/4)=Z/2$for i=1,2, and if I look at the reduced cohomology long exact sequence $H^0(M,pt,Z/2)=H^0(M,pt,z/4)=0$. Then the Bockstein from $H^2$to $H^1$will not be trivial, otherwise we will have $0\to Z/2\to Z/2\to Z/2\to 0$ contradiction. Although I used the reduced cohomology, since I only use maps to K(z/2,2) and K(Z/2,1) to represent $H^1$ and $H^2$ I am fine. From this example, we get a non null-homotopic map from $CP^{\infty}$ to $RP^\infty$! – Ying Zhang Apr 4 '10 at 4:05
Sorry about my last sentence, from the above sequence I got a map from $RP^{\infty}\to K(Z/2,2)$. To get a map from $RP^{\infty}\to CP^{\infty}$, we should look at the coefficient sequence $0\to Z\to Z\to Z/2Z\to 0$ which does the job. Geometrically it is not too hard to find an interesting map from $RP^{\infty}$ to $CP^{\infty}$ just by quotient out more things. – Ying Zhang Apr 4 '10 at 18:15
For a more explicit example than Chris's, consider the map from the (2-dimensional) torus to a sphere that collapses the 1-skeleton of the usual CW complex and takes the 2-cell to the 2-cell of the sphere. The torus is $K(\mathbb{Z}^2,1)$, so this necessarily gives zero maps on homotopy, but it's also pretty clearly not null-homotopic.
-
Nice! So this will give us examples when X and Y are both manifolds:) – Ying Zhang Apr 4 '10 at 4:11
You can also generalize this to the case where X is an $n$-manifold and a $K(G,1)$ (or even just having nonzero homotopy groups only up to $n-1$), and $X \to S^n$ is the map that collapses everything outside a Euclidean ball to a point. – Kevin Casto May 27 at 23:26
Even if you ask that $f$ induces trivial maps on all (singular) homology and cohomology groups, there are still easy manifold examples. (This actually arises as an exercise in Hatcher's AT).
For instance, let $f:T^3\rightarrow S^2$ be the composition $T^3\rightarrow S^3\rightarrow S^2$, where the map from $T^3$ to $S^3$ is simply collapsing the 2-skeleton to a point, and the map from $S^3$ to $S^2$ is the Hopf map.
As others have mentioned, since $T^3$ is a $K(\mathbb{Z}^3, 1)$, if follows that $f$ induces trivial maps on homotopy groups.
Since the Hopf map induces trivial maps on homology and cohomology, it follows that $f$ does as well.
Finally, to see that $f$ is NOT nullhomotopic, assume it is. Since the map from $S^3$ to $S^2$ is a fiber bundle, it has the homotopy lifting property. Hence, we can lift the homotopy of $f$ to a homotopy $G:I\times T^3\rightarrow S^3$ where $G_0$ is the above map from $T^3$ to $S^3$ and $G_1$ is is a map from $T^3$ to $S^1\subseteq S^3$, the preimage of a point in $S^2$ under the Hopf map.
But $G_0$ has degree 1, while $G_1$ has degree 0, a contradiction.
-
Thanks. Btw, the exercise you mentioned is on Hatcher p.392 problem 34. – Ying Zhang Apr 4 '10 at 18:18
I asked a very similar question a few months ago, and got some excellent answers.
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https://www.physicsforums.com/threads/question-about-eigenvectors.101776/ | 1. Nov 28, 2005
### jd1828
I am studying for a test that I have today and just need a quick anwser. the problem is find the eigenvector of [8 -10;2 -1] where ; means skip to a new line. So I get lamda=4,3 then try to find eigenvector for lamda=4. I get the vector [5/2;1] which works but the anwser in the book is [2;5]. Both anwsers seem to be right but is there some rule to say that the book anwser is better.
2. Nov 28, 2005
### TD
Eigenvectors are always determined up to an arbitrary (non-zero) multiplication factor. The book apparently preferred not to have fractions but the smallest possible integers, which is usually done. Both are fine though
3. Nov 28, 2005
### matt grime
its answer doesn't have any fractions in it, that is considered better, just as you should replace 1/sqrt(2) with sqrt(2)/2. Its answer is also valid even if we only consider vectors with entries in the integers (which we may want to) so it is more general. Your answer is fine if we remember that we CAN divide by two, but it holds in greater generality than that.
Just remember that if you can you should clear denominators in an answer, and since e-vectors are only defined up to scalar multiples why not? PERSONAL OPINION and NOT in any way authoritative is that I wouldn't mark you down for your answer. But that is purely MY opinion and is not necessarily that of your teacher/marker.
actually, have just noticed your answer and the books are not compatible, bnut that is probably a typo
(5/2,1) is not a scalar multiple of (2,5)
(5/2,1) and (5,2) yes or
(2/5,1) and (2,5) also ok
4. Nov 28, 2005
### jd1828
Thanks a lot! Ill stick to using the smallest possible integers since thats what the book does and I wouldnt want to loose points on something so easy. The anwser in the book was [5,2] I just typed it wrong.
5. Nov 28, 2005
### HallsofIvy
As pointed out before, if v is an eigenvector corresponding to eigenvalue $\lambda$ the so is any multiple of v. Yes, [5/2, 1] is an eigenvector corresponding to eigenvalue 4. Multiplying that vector by 2, so is [5, 2]. Some people prefer to have integer components (which is not always possible), some people prefer to write eigenvectors so that have unit length. Since [5,2] has length $\sqrt{29}$ an unit eigenvector would be $$\left[\frac{5}{\sqrt{29}},\frac{2}{\sqrt{29}}\right]$$
6. Nov 29, 2005
### jd1828
Why would you want to use the unit vector. We never really went over that in class but there is a question on a lab that uses it. Where we had to show that the vectors are mutually orthogonal. I dont have the exact question with me but Im still working on the anwser.
7. Nov 30, 2005
### HallsofIvy
If you have a basis for your vector space where the vectors in the basis are "orthonormal"- orthogonal and have length 1- then calculations are much simplified. If, in addition, your basis vectors are all eigenvectors of some linear transformation, then working with that linear transformation is very simple.
8. Nov 30, 2005
### jd1828
Im home now so I can give you the exact question.
Given matrix a=[3,3,2; 3,6,5; 2,5,11]
a) determine eigenvalues.
This part I can do just fine
b) Determine unit eigenvectors with magnitude one.
So I just take each eigenvector and divide by the lenght?
c) Show that the unit eigenvectors are mutually orthogonal.
This I have no idea what to do. Ive been reading through the book and it looks like I have to show that two unit vectors mutiplied together are zero.
9. Nov 30, 2005
### TD
Indeed, just as HallsofIvy illustrated.
Yes, if you mean the scalair product (inner product).
10. Nov 30, 2005
### HallsofIvy
Good! That's the hard part!
Yes, that's exactly right.
Reading the book is always a good idea! Two vectors are orthogonal if and only if their dot product is equal to 0. Becareful to say "dot" product- there are several different kinds of multiplication defined for vectors. Take the dot product of each pair of vectors and see what you get.
11. Nov 30, 2005
### jd1828
Thanks a lot! I tried it out and the dot products do equal zero. | 2018-02-18T20:35:14 | {
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http://rozmierka.eu/journal/8c5079-disjoint-vs-independent | In mutually exclusive events, the occurrence of one event will result in the non-occurrence of the other. Independent means if one happens it doesn't affect whether or not the other happens. Does this use of the perfect actually express something about the future? JavaScript is disabled. You stand at the basketball free-throw line and make 30 attempts at at making a basket. If possible, please give more than one example and counterexample. Disjoint means the two events are mutually exclusive -- if one happens than the other can't happen. The probability of the empty set is zero; therefore, the event that both G and H occur is impossible.This means that G and H are mutually exclusive.. How to Identify Independent Events. Disjoint means the two events are mutually exclusive -- if one happens than the other can't happen. Ahhh i seee, so if disjoint, they cannot both occur, so the intersection is the empty set. That is the reason you are having trouble finding an example of this case; it is a pathological case that is not important. Quick way to move an object some distance from one external vertex to another external vertex? I wouldn't say that fair coin tosses are in any way dependent on each other. You are right though: in the more general case beyond a binary partition, you could have multiple disjoint events that don't cover the sample space and each have probability zero. Empty set disjoint with itself paradoxical? You and a randomly selected student from your class both earn A's in this course. Is there an example of two causally dependent events being logically (probabilistically) independent? Well, tossing a fair coin leads to series of events that are both independent and disjoint. For independence, # P(A \cap B) = P(A) * P(B) #, If two disjoint events are independent then #P(A \cap B) = 0 = P(A) * P(B) #. rev 2020.11.2.37934, Sorry, we no longer support Internet Explorer, The best answers are voted up and rise to the top, Cross Validated works best with JavaScript enabled, Start here for a quick overview of the site, Detailed answers to any questions you might have, Discuss the workings and policies of this site, Learn more about Stack Overflow the company, Learn more about hiring developers or posting ads with us. Disjoint events are mutually exclusive, which is a strong form of statistical dependence (so if you know event A occurred you know that B definitely did not occur and vice versa), meaning. (C64). It is more likely that rates of milk-consumption differ for Trump-voters and non-Trump voters, for all sorts of indirect reasons. To learn more, see our tips on writing great answers. In particular, within the population of US adults, males drink more milk, and are also more likely to vote for Trump, so if this is one of the main indirect effects, I would expect that there would be some small positive correlation between milk-consumption and voting for Trump. Ok, what about the union, i.e what aboue p(A U B ) if disjoint or independent, are we allowed to sum for both? Intro to probability 2: Independent and disjoint ( - YouTube Specific Multiplication Rule. Suggestions for braking with severe osteoarthritis in both hands, Why does the VIC-II duplicate its registers? What's some examples of the fact disjoint events are dependent? Disjoint vs. independent. Two events are mutually exclusive if they can't both happen. When you think about independent events, you should be thinking about the probabilities of the events; the likelihood of one event happening is not influencing the likelihood of another. 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Two events A and B are said to be independent if the outcome of event A doesn’t affect the outcome of event B and vice versa. Disjoint/independent: Disjoint events are only statistically independent if each outcome has probability zero or one (which is a trivial case). They are still disjoint technically because their intersection is still empty. Conversely, in independent events, occurrence of one event will have no influence on the occurrence of the other. Detailed explanation: what is "dayspring"? As a rough rule of thumb (which is slightly exaggerated, but not much) everything is correlated with everything. Disjoint implies dependent. Although they are unlikely to be causally related, it would surprise me if voting-preference is uncorrelated with milk-consumption. Cross Validated is a question and answer site for people interested in statistics, machine learning, data analysis, data mining, and data visualization. How can a hive mind secretly monetize its special ability to make lots of money? So, in Part 1, red and orange in *one* draw are clearly disjoint. So when we have evens and we say the probability of their union is equal to their sum, does this mean they are independent or disjoint ? Why sister [nouns] and not brother [nouns]? An event is deemed independent if it offers no information about other events. Why does a blocking 1/1 creature with double strike kill a 3/2 creature? Making statements based on opinion; back them up with references or personal experience. Only valid for independent events P(A and B) = P(A) * P(B) Example 3: P(A) = 0.20, P(B) = 0.70, A and B are independent. Stack Exchange network consists of 176 Q&A communities including Stack Overflow, the largest, most trusted online community for developers to learn, share their knowledge, and build their careers. Adjective agreement-seems not to follow normal rules. Set Theory, Logic, Probability, Statistics, Genomic data 'catches corals in the act' of speciation and adaptation, Bats can predict the future, researchers discover, Neanderthal children grew and were weaned similarly to modern humans. Are these definitions correct? The Addition Rule for Disjoint Events (Rule Four) The General Addition Rule for which the events need not be disjoint (Rule Five) In order to complete our set of rules, we still require two Multiplication Rules for finding P(A and B) and the important concepts of independent events and conditional probability. I'm trying to understand joint probabilities and dependence/independence. The last 3 times you went to the dentist for your 6-month checkup, it rained as you drove to her... You roll a balanced die two times. Your other examples look correct (at least in spirit, I have not surveyed voter milk drinking habits). Use MathJax to format equations. What is the difference between disjoint and independent events, how will the 2 affect calculations involving them? Independent means if one happens it doesn't affect whether or not the other happens. Are these correct? Asking for help, clarification, or responding to other answers. If you draw one card from a deck of cards, what is the probability that it is a heart or a diamond? around the world. Joint/disjoint and dependent/independent events, rates of milk-consumption among US adults differ by age and sex, voting patterns also differ by age and sex, Creating new Help Center documents for Review queues: Project overview, 2020 Moderator Election Q&A - Questionnaire, 2020 Community Moderator Election Results. In your first 350 flips, you have obtained 300 tails and 50 heads. Is the nucleus smaller than the electron? Read on to find out more about dependent events vs independent events. For example, rates of milk-consumption among US adults differ by age and sex, and voting patterns also differ by age and sex. For a better experience, please enable JavaScript in your browser before proceeding.
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https://math.stackexchange.com/questions/1267423/understanding-a-proof-showing-that-for-any-prime-p-there-are-integers-x-and | Understanding a proof showing that for any prime $p$, there are integers $x$ and $y$ such that $p|(x^2+y^2+1)$.
I asked this question a couple days ago: Show that for any prime $p$, there are integers $x$ and $y$ such that $p|(x^{2} + y^{2} + 1)$. But I asked it as a guest, and I could not comment on the answers I got as I did not have enough reputation.
I have found a proof from Warwick University that will help me answer this question, but there is one sentence I am unsure of. If anybody could help me out on it, that would be amazing!
Here it is:
There are $\frac{p+1}{2}$ numbers of the form $x^2 \pmod{\!p}$. There are $\frac{p+1}{2}$ numbers of the form $-1-y^2 \pmod{\!p}$ (because to get the numbers of the form $-1-y^2$, multiply the squares by $-1$ and subtract $1$… this will not change how many there are). There are exactly $p$ numbers mod $p$. Note that $\frac{p+1}{2}$ is more than half of $p$. So the set of numbers of the form $x^2$ and those of the form $\frac{p+1}{2}$ must have at least one common element (if not, we will have too many distinct numbers mod $p$). So there are some $x$ and $y$ such that $x^2\equiv -1-y^2 \pmod{\!p}$. Thus $x^2 + y^2 +1\equiv 0 \pmod{\!p}$.
I understand it up until the bolded sentence. Did they mean to write numbers of the form $\frac{p+1}{2}$? Or did they instead mean numbers of the form $-1-y^2$?
No matter what they meant, would anybody be able to explain more how having at least one overlapping element in both sets tells you that solutions to that congruence exists? How did they come to that conclusion that there is a common element? What would happen if there wasn't a common element? Thank you!
• I think that they meant "of the form $-1-y^2$" and this is axactly what I wrote in my answer for your question – Elaqqad May 8 '15 at 14:36
As you have guessed, there is an error (maybe just a typo) in that sentence. It should be
So the set of numbers of the form $x^2$ and those of the form $-1-y^2$ must have at least one common element (if not, we will have too many distinct numbers mod $p$).
The point is that if you give all possible values to $x$ modulo $p$, all possible values to $y$ and then work out all the values of $x^2$ and of $-1-y^2$ modulo $p$, you will get $\frac{p+1}2$ different values for $x^2$ and $\frac{p+1}2$ different values for $-1-y^2$, that's $p+1$ answers altogether.
But you cannot get $p+1$ different answers as there are only $p$ possibilities modulo $p$. So two of your answers must be the same: either
• two of the $x^2$ values are the same; or
• two of the $-1-y^2$ values are the same; or
• an $x^2$ value is the same as a $-1-y^2$ value.
But we have already noted that the first two of these options are impossible, so the third one must be true.
If it helps, look at a specific example, say $p=7$. There are four possible values of $x^2$, namely $$0,\ 1,\ 4,\ 2.$$ Now without working them out, you can see that there will be four values of $-1-y^2$: so we will have eight numbers written down: so two of them must coincide. To confirm this, the values of $-1-y^2$ are $$-1,\ -2,\ -5,\ -3,$$ that is, $$6,\ 5, 2,\ 4,$$ and you can see that the result $2$ occurs twice. (In fact, so does the result $4$, giving another possible solution.)
Exercise: show that the "other possible solution" is really the same solution.
• You've saved my life, thank you so much! Completely understand now! – user236043 May 5 '15 at 4:52
A more general version of this result can be found in Herstein's : Topics in Algebra text. Kindly see Lemma 7.1.7 Page $359$. The result states that
• If $F$ is a finite fied and $\alpha \neq 0$ and $\beta \neq 0$ are two elements of $F$ then we can find elements $a$ and $b$ in $F$ such that $1 + a \cdot\alpha^{2} + b \cdot \beta^{2}=0$.
• This is not really a generalisation because this says only that $1,\alpha^2,\beta^2$ are dependant which is not really the same thing as saying that there exists some $x,y$ such that $1+x^2+y^2=0$ (it's correct but what I mean is that your answer is not a generalisation it's another separte result ) – Elaqqad May 8 '15 at 14:35 | 2019-08-19T11:51:04 | {
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https://math.stackexchange.com/questions/2337650/finding-the-basis-for-the-null-space-of-4-times-4-matrix | # Finding the basis for the null space of $4\times 4$ matrix
My linear algebra professor gave us a practice worksheet for our upcoming exam, and the answer key he gave has me very confused. I'm not entirely sure his answers are correct, which is why I thought i'd bring the question to you good folks. The question is as follows:
Find a basis for the solution space of Ax = 0 if
$$A = \begin{bmatrix} 1&1&0&0\\ -2&-2&0&0\\ 0&0&1&-1\\ -1&-1&0&1 \end{bmatrix}$$
Reduced row echelon form yields:
$$\left[ \begin{array}{cccc|c} 1&1&0&0&0\\ 0&0&1&0&0\\ 0&0&0&1&0\\ 0&0&0&0&0 \end{array} \right]$$
$\implies$
$x_1 = -x_2$
$x_3 = 0$
$x_4 = 0$
$\implies$ Basis for null(A) = $\begin{bmatrix} -1\\ 1\\0\\0 \end{bmatrix}$
His solution, however, gives:
Basis for null(A) = $\begin{bmatrix} -1\\ 1\\0\\0 \end{bmatrix}$ , $\begin{bmatrix} 0\\ 0\\-1\\1 \end{bmatrix}$
Did I make a mistake here? Or is this answer incorrect. I have always always been under the assumption that $dim(\mathbb R^4) = rank(A) + nullity(A)$ per the rank-nullity theorem. In this case:
$4 = 3 + nullity(A) \implies nullity(A) = 1$ meaning, to me, the null space should contain 1 vector and not 2.
If i'm totally off base here please let me know, as I really want to understand this stuff. Thanks!
• Did you check your lecturer's solution? Take the vectors $v_1$ and $v_2$ he gave you and see whether $Av_1$ and $Av_2$ are both the zero vector? – Lord Shark the Unknown Jun 27 '17 at 4:43
• It did not give the zero vector for A$v_2$, although I was not aware of that check, thank you for sharing. So i'm guessing this means his solution is wrong? – FuegoJohnson Jun 27 '17 at 4:49
• Yes, this means his solution is wrong! – Lord Shark the Unknown Jun 27 '17 at 4:55
• @FuegoJohnson This sanity check is nothing more than applying the definition of a null space to the purported basis of a null space. – amd Jun 27 '17 at 6:58
• Quite likely there’s a typo in the matrix $A$. A small change to it will make the null space two-dimensional. – amd Jun 27 '17 at 7:03
• Since the eigenspace is the null space of $det(A - \lambda I)$, it follows that the dimension of the eigenspace is $n - rank(A - \lambda I)$ for any square matrix. Is this correct? – FuegoJohnson Jun 27 '17 at 5:56 | 2019-11-19T21:13:39 | {
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https://stats.stackexchange.com/questions/318098/how-should-i-formulate-this-hypothesis-test/318099 | How should I formulate this hypothesis test?
I have the following problem solved by my professor:
A company of high resistance painting products claims that, at the most, 1% of its containers have a weight outside specification limits. To investigate the plausibility of this assertion, a sample of n = 45 containers has been taken, which is such that only one container is defective. At the level α = 0.05, is there evidence against the manufacturer’s claim?
The way I thought I could test this is by setting: $$H_0:p=0.01$$ $$H_1:p<0.01$$, and testing this, I get that the null hypothesis is not rejected so that there is evidence against the manufacturer's claim. However, my professor solved it by testing: $$H_0:p=0.01$$ $$H_1:p>0.01$$ , and concludes that there is not enough evidence. I am really confused. How should this be approached?
Notice that the test you want to claim is that $p=0.01$ at most. So you need to formulate the null hypotheses as what the company (or whoever makes a claim, depending on the problem statement) claims, and the alternative as something that would prove them (the company) wrong.
In this case, they claim "at most" p=0.01, so the alternative, what would prove them wrong is $p>0.01$, so $p$ is more than 1%.
• Thank you for your reply. I understand what you are saying, and I will start solving my problems following your procedure. Would my interpretation be wrong then? Intuitively, both tests should reach the same conclusion right? – Bee Dec 10 '17 at 22:08
• No, both tests will not reach the same conclusion. When you solve hypotheses testing problems, the statement says something like "person/entity ABC claims that object XYZ is this way". You have to think in terms of what would prove their claim to be wrong. In this case, if they say "at most 1% defects", showing that it is less than 1% ($p<0.01$) would NOT prove them wrong, it would in fact prove them even more right. The opposite of "at most 1%" is "more than 1%", this is your alternative hypothesis. – Anna SdTC Dec 10 '17 at 22:14
If p < 0.01 that would mean that the maximum proportion of defects is even less that 1%, so the alternative hypothesis set by your professor seems right. Your alternative hypothesis doesn't really play a role if you're only concerned with Type I error, that is, the probability that you will reject a true null hypothesis, which in your case is set as 5% ($\alpha$ being 0.05).
And notice that your $H_0$ not getting rejected means that there is more that 5% chance that the data was indeed sampled from a binomial distribution with p = 0.01 (where success is defined as a defect), which I guess is how you model your null hypothesis distribution. So in fact you too should conclude that there is not enough evidence to decide whether manufacturer is telling the truth or not.
In my conservative approach to this problem, the company's claim should be set as the alternative hypothesis, not the null hypothesis. It is incumbent upon the company to prove their quality. By setting the company's claim as the null hypothesis, their claim is supported up and until enough evidence is accumulated to reject their claim. This is an approach to the problem, but effectively creates a logic that the company's quality is assumed good until proven otherwise. We usually assume a world in which claims require evidence, not one in which claims are assumed true. In addition, the hypothesis tests should cover the entire range of possible results giving:
$$H_0:p(defects) > 0.01 \\ H_a:p(defects) \leq 0.01$$
The test statistic is 1/45 (2.2%) with a standard error of 2.2%. The difference between 1% claimed and 2.2% measured is 1.2%. Standardizing, we obtain a z-score of 0.54 and a p-value of 0.71. Based on the sampled data, we cannot reject the null hypothesis that the probability of defects is >0.01 at a one-sided alpha $\le$ 0.05. Thus, I arrive at the same answer, but protect myself against accepting the company's claim as the result of insufficient data.
• This analysis is the opposite of "conservative." To see why, suppose that all 45 objects in the sample were defective. Your test statistic of $100\%$ would have an enormous Z-score and you would obtain a p-value of $1.00$, thereby not concluding that the company's claim is false--in violation of all common sense and the evidence. – whuber Dec 11 '17 at 17:10
• While defective, if I'm asking a one-sided question related to the left tail, am I not still internally consistent with my setup regardless of numerous failures? – Todd D Dec 11 '17 at 17:24 | 2019-12-13T09:58:45 | {
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https://math.stackexchange.com/questions/978385/a-1-cap-a-2-cap-cdots-cap-a-n-ne-emptyset-holds-for-all-n-must-it-be | # $A_1 \cap A_2 \cap \cdots \cap A_n \ne \emptyset$ holds for all $n$. Must it be that $\bigcap_{n = 1}^{\infty} A_n \ne \emptyset$?
Let $A_1, A_2, A_3, \,\ldots$ be sets such that $A_1 \cap A_2 \cap \cdots \cap A_n \ne \emptyset$ holds for all $n$.
Must it be that $\bigcap_{n = 1}^{\infty}A_n \ne \emptyset$?
I answered no. Here is my "proof".
Define $A_n = \{n+1, n+2, \,\ldots\}$.
Then $A_1 \cap A_2 \cap \cdots \cap A_n = \{n+1, n+2, \,\ldots\}$.
For all integers $n$, $\{n+1, n+2, \,\ldots\} \ne \emptyset$, since it contains $n+1$, and there is no largest integer.
Now consider $\bigcap_{n = 1}^{\infty} A_n$.
Suppose it contains an integer $m$. But $m \notin A_{m}$ by definition, hence $m \notin \bigcap_{n = 1}^{\infty} A_n$, so $\bigcap_{n = 1}^{\infty} A_n = \emptyset$.$$\tag*{\blacksquare}$$
Is that correct?
I ask because in my lecture notes, the definition of $\bigcap_{n = 1}^{\infty} A_n$ is $\{x : x \in A_n \ \forall n \}$ which (to me) seems to be equivalent to "$A_1 \cap A_2 \cap \cdots \cap A_n \ne \emptyset$ holds for all $n$", which would of course mean the answer is yes.
• no is the correct answer- and your example works. – voldemort Oct 17 '14 at 15:32
• Looks fine to me. +1 for showing your work. – Timbuc Oct 17 '14 at 15:33
• No is the correct answer, your example works, and this is still compatible with the definition you were given. $\bigcap_{n=1}^\infty A_n \neq \emptyset$ does imply that for any $m$ you have $\bigcap_{n=1}^m A_n \neq \emptyset$. But it is stronger, as your example demonstrates. – Ian Oct 17 '14 at 15:34
• I think your confusion comes down to order of quantifiers. The finite intersections all being nonempty means $(\forall m \in \mathbb{N}) (\exists x)(\forall n \in \{ 1,\dots,m \}) \, x \in A_n$. This means $x$ can depend on $m$ because $x$ is bound under the scope of the $m$ quantifier. But the countable intersection being nonempty means $(\exists x) (\forall n \in \mathbb{N}) \, x \in A_n$. Here $x$ is independent of $n$. – Ian Oct 17 '14 at 15:40
• Another interesting example, not discrete, is to take $A_n = (0,\frac1n)$. Then $$A_1\cap\cdots\cap A_n = A_n$$ but there is no point in $\bigcap_{n=1}^{\infty}A_n$ because any such point $a$ would have to satisfy $0<a<\frac1n$ for all $n$, which is impossible. – MPW Oct 17 '14 at 15:41
The key here, as Ian's comment says, is that the first statement $\;\langle \forall n :: \langle \cap i : i \leq n : A_i \rangle \not= \emptyset \rangle\;$ is (by the definitions) equivalent to $$\langle \forall n :: \langle \exists x :: \langle \forall i : i \leq n : x \in A_i \rangle \rangle \rangle$$ while the second statement $\;\langle \cap i :: A_i \rangle \not= \emptyset\;$ is equivalent to $$\langle \exists x :: \langle \forall i :: x \in A_i \rangle \rangle$$
These two are not equivalent, as you correctly show by your counterexample of $\;A_i = \{j | j > i\}\;$ which makes the first statement true (by choosing $\;x := n+1\;$) and the second false.
$$\calc \langle \forall n :: \langle \exists x :: \langle \forall i : i \leq n : x \in A_i \rangle \rangle \rangle \calcop{\Leftarrow}{logic: strengthen by weakening range of \;\forall i\;} \langle \forall n :: \langle \exists x :: \langle \forall i :: x \in A_i \rangle \rangle \rangle \calcop{\Leftarrow}{logic: remove unused \;\forall n\;} \langle \exists x :: \langle \forall i :: x \in A_i \rangle \rangle \endcalc$$
Take $A_n = \left(0, \dfrac{1}{n}\right)$, then $\displaystyle \bigcap_{k=1}^n A_k = \left(0,\dfrac{1}{n}\right)$, and $\displaystyle \bigcap_{k=1}^\infty A_k = \emptyset$.
• Did you read the question...? – Najib Idrissi Oct 17 '14 at 18:11
• @NajibIdrissi: yes he did,this is a counter example. – Arjang Oct 18 '14 at 10:08
• @Arjang The question isn't asking for a counterexample, the question is asking to check the proof. Notice that the question already contains a counterexample. – Najib Idrissi Oct 18 '14 at 10:15
• @NajibIdrissi : Ahaa! I didn't read the question either :) – Arjang Oct 18 '14 at 10:23 | 2019-10-13T23:35:47 | {
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https://tech.io/playgrounds/257/greatest-common-divisor | greatest common divisor
cehsu
104 views
Greatest Common Divisor
In mathematics, the greatest common divisor (gcd) of two or more integers, which are not all zero, is the largest positive integer that divides each of the integers. For example, the gcd of 8 and 12 is 4.
The greatest common divisor is also known as the greatest common factor (gcf), highest common factor (hcf), greatest common measure (gcm), or highest common divisor.
Check for the gcd
Write a getGCD function that takes two numbers as parameters and returns the gcd.
Sample Input and Output
Input: 8 12 Output: 4
Optimization
Check out Euclid's algorithm and see how much things speed up. How can you account for this in terms of time complexity?
See:
https://en.wikipedia.org/wiki/Euclidean_algorithm#Worked_example
https://en.wikipedia.org/wiki/Greatest_common_divisor#Complexity_of_Euclidean_method
Write the getGCD function
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function getGCD (numOne, numTwo) {
}
module.exports = getGCD;
XXXXXXXXXXXXXXXXXXXXXXXXXXXXXXXXXXXXXXXXXXXXXXXXXX
Open Source Your Knowledge: become a Contributor and help others learn. | 2017-08-20T02:18:34 | {
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https://ocw.mit.edu/courses/aeronautics-and-astronautics/16-90-computational-methods-in-aerospace-engineering-spring-2014/numerical-methods-for-partial-differential-equations/eigenvalue-stability-of-finite-difference-methods/1690r-matrix-stability-for-finite-difference-methods/ | # 2.7 Eigenvalue Stability of Finite Difference Methods
## 2.7.2 Matrix Stability for Finite Difference Methods
As we saw in Section 2.3.2, finite difference (or finite volume) approximations can potentially be written in a semi-discrete form as,
$\frac{{\rm d}U}{{\rm d}t} = AU + b. \label{equ:semidiscrete}$ (2.133)
While there are some PDE discretization methods that cannot be written in that form, the majority can be. So, we will take the semi-discrete Equation (2.133) as our starting point. Note: the term semi-discrete is used to signify that the PDE has only been discretized in space.
Let $$U(t)$$ be the exact solution to the semi-discrete equation. Then, consider perturbation $$e(t)$$ to the exact solution such that the perturbed solution, $$V(t)$$, is:
$V(t) = U(t) + e(t).$ (2.134)
The questions that we wish to resolve are: (1) can the perturbation $$e(t)$$ grow in time for the semi-discrete problem, and (2) what the stability limits are on the timestep for a chosen time integration method.
First, we substitute $$V(t)$$ into Equation (2.133),
$$\displaystyle \frac{{\rm d}V}{{\rm d}t}$$ $$\displaystyle =$$ $$\displaystyle AV + b$$ (2.135) $$\displaystyle \frac{{\rm d}(U+e)}{{\rm d}t}$$ $$\displaystyle =$$ $$\displaystyle A(U+e) + b$$ (2.136) $$\displaystyle \frac{{\rm d}e}{{\rm d}t}$$ $$\displaystyle =$$ $$\displaystyle A e.$$ (2.137)
Thus, the perturbation must satisfy the homogeneous equation, $$e_ t = A e$$. Having studied the behavior of linear system of equations in Section 1.6.2, we know that $$e(t)$$ will grow unbounded as $$t \rightarrow \infty$$ if any of the real parts of the eigenvalues of $$A$$ are positive.
The problem is that determining the eigenvalues of $$A$$ can be non-trivial. In fact, for a general problem finding the eigenvalues of $$A$$ can be about as hard as solving the specific problem. So, while the matrix stability method is quite general, it can also require a lot of time to perform. Still, the matrix stability method is an indispensible part of the numerical analysis toolkit.
As we saw in the eigenvalue analysis of ODE integration methods, the integration method must be stable for all eigenvalues of the given problem. One manner that we can determine whether the integrator is stable is by plotting the eigenvalues scaled by the timestep in the complex $$\lambda {\Delta t}$$ plane and overlaying the stability region for the desired ODE integrator. Then, $${\Delta t}$$ can be adjusted to attempt to bring all eigenvalues into the stability region for the desired ODE integrator.
## Matrix Stability of FTCS for 1-D convection
Earlier, we used a forward time, central space (FTCS) discretization for 1-d convection,
$\frac{U_ i^{n+1}-U_ i^ n}{{\Delta t}} + u_{i}^ n\delta _{2x} U^ n_{i} = 0. \label{equ:con1d_ ftcs}$ (2.138)
Since this method is explicit, the matrix $$A$$ does not need to be constructed directly, rather Equation (2.138) can be used to find the new values of $$U$$ at each point $$i$$. However, if we are interested in calculating the eigenvalues to analyze the eigenvalue stability, then the $$A$$ matrix is required. The following script does exactly that (i.e. calculates $$A$$, determines the eigenvalues of $$A$$, and then plots the eigenvalues scaled by $${\Delta t}$$ overlayed with the forward Euler stability region). The script can set either the inflow/outflow boundary conditions described in Example 2.3.3, or can set periodic boundary conditions. We will look at the eigenvalues of both cases.
% This MATLAB script calculates the eigenvalues of
% the one-dimensional convection equation discretized by
% finite differences. The discretization uses central
% differences in space and forward Euler in time.
%
% Periodic bcs are set if periodic_flag == 1.
%
% Otherwise, an inflow (dirichlet) bc is set and at
% the outflow a one-sided (backwards) difference is used.
%
clear all;
periodic_flag = 1;
% Set-up grid
xL = -4;
xR = 4;
Nx = 21; % number of points
x = linspace(xL,xR,Nx);
% Calculate cell size in control volumes (assumed equal)
dx = x(2) - x(1);
% Set velocity
u = 1;
% Set timestep
CFL = 1;
dt = CFL*dx/abs(u);
% Set bc state at left (assumes u>0)
UL = exp(-xL^2);
% Allocate matrix to hold stiffness matrix (A).
%
A = zeros(Nx-1,Nx-1);
% Construct A except for first and last row
for i = 2:Nx-2,
A(i,i-1) = u/(2*dx);
A(i,i+1) = -u/(2*dx);
end
if (periodic_flag == 1), % Periodic bcs
A(1 ,2 ) = -u/(2*dx);
A(1 ,Nx-1) = u/(2*dx);
A(Nx-1,1 ) = -u/(2*dx);
A(Nx-1,Nx-2) = u/(2*dx);
else % non-periodic bc's
% At the first interior node, the i-1 value is known (UL).
% So, only the i+1 location needs to be set in A.
A(1,2) = -u/(2*dx);
% Outflow boundary uses backward difference
A(Nx-1,Nx-2) = u/dx;
A(Nx-1,Nx-1) = -u/dx;
end
% Calculate eigenvalues of A
lambda = eig(A);
% Plot lambda*dt
plot(lambda*dt,'*');
xlabel('Real \lambda\Delta t');
ylabel('Imag \lambda\Delta t');
% Overlay Forward Euler stability region
th = linspace(0,2*pi,101);
hold on;
plot(-1 + sin(th),cos(th));
hold off;
axis('equal');
grid on;
Figures 2.21 and 2.22 show plots of $$\lambda {\Delta t}$$ for a CFL set to one. Recall that for this one-dimensional problem, the CFL number was defined as,
$\mathrm{CFL} = \frac{|u|{\Delta t}}{{\scriptstyle \Delta } x}.$ (2.139)
In the inflow/outflow boundary condition case (shown in Figure 2.21) the eigenvalues lay slightly inside the negative real half-plane. As they move away from the origin, they approach the imaginary axis at $$\pm i$$. The periodic boundary conditions give purely imaginary eigenvalues but these also approach $$\pm i$$ as the move away from the origin. Note that the periodic boundary conditions actually give a zero eigenvalue so that the matrix $$A$$ is actually singular (Why is this?). Regardless what we see is that for a $$\mathrm{CFL} = 1$$, some $$\lambda {\Delta t}$$ exist which are outside of the forward Euler stability region. We could try to lower the timestep to bring all of the $$\lambda {\Delta t}$$ into the stability region, however that will prove to be practically impossible since the extreme eigenvalues approach $$\pm \alpha i$$ (i.e. they are purely imaginary). Thus, no finite value of $${\Delta t}$$ exists for which these eigenvalues can be brought inside the circular stability region of the forward Euler method (i.e. the FTCS is unstable for convection).
Figure 2.21: $$\Lambda {\Delta t}$$ distribution for one-dimensional convection example with Dirichlet boundary conditions. Note $${\Delta t}$$ set such that $$CFL=1$$.
Figure 2.22: $$\Lambda {\Delta t}$$ distribution for one-dimensional convection example with periodic boundary conditions. Note $${\Delta t}$$ set such that $$CFL=1$$.
We may also be interested in what happens to the eigenvalue spectrum of $$A$$ when we change $$\Delta x$$ keeping the CFL constant. Figures 2.23 and 2.24 plot the eigenvalue spectrum for both periodic and Dirichlet boundary conditions refining $$\Delta x$$ by a factor of 10. For the periodic BC case, again we see purely imaginary eigenvalues approaching $$\pm i$$ (though, of course, we have more eigenvalues as $$A$$ is now a larger system). For the Dirichlet case, the eigenvalues again lie slightly inside the negative real half-plane, though in this case closer to the imaginary axis than for the coarser system. In fact, the eigenvalues of the Dirichlet problem approach those of the periodic problem in the limit as $$\Delta x \to 0$$.
Figure 2.23: $$\Lambda {\Delta t}$$ distribution for one-dimensional convection example with Dirichlet boundary conditions. Note $$\Delta x=\frac{1}{20}$$ and $${\Delta t}$$ set such that $$CFL=1$$.
Figure 2.24: $$\Lambda {\Delta t}$$ distribution for one-dimensional convection example with periodic boundary conditions. Note $$\Delta x = \frac{1}{20}$$ and $${\Delta t}$$ set such that $$CFL=1$$. | 2021-09-20T16:46:49 | {
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http://prestavamkourit.cz/bunk-beds-crt/23d4a9-do-angles-in-a-kite-add-up-to-360 | All of the exterior angles of a polygon add up to 360°. Parallel lines … A kite has two pairs of equal sides. 3. Angles around a point add up to what value? If you were to walk around a four sided shape until you reach the start you will have rotated by 360 degrees. x + y + y + 16 = 360. x + y + y = 344 <----(360 -16) Now, let me explain this. The angles of a quadrilateral add up to 360°. Yes, because all the angles in any quadrilateral (4-sided polygon) add up to 360. One pair of opposite angles are equal in a kite and its 4 interior angles add up to 360 degrees. Angles in a triangle add up to 180° and in quadrilaterals add up to 360°. Abhinav. We have the following theorem about the angles of a quadrilateral. 1 Questions & Answers Place. so we form the following equation. The angles at a point (full turn) add up to 360 degrees. Each quadrant of the unit circle is composed of many triangles that make up the quadrant. The angles inside any shape with 5 sides add up … To show that the angle sum of a triangle equals 180 degrees, draw a triangle, tear the angles and rearrange them into a straight line.Remember that the number of degrees in a straight line is 180 degrees. Square. Which type of angle is between 180° and 360°? Therefore we now have two triangles within one Quadrilateral. The angles in a parallelogram will always add up to 360°. Kite; That does not seem to many does it. The angle opposite the 30° is 180° - 30° = 150°. Yes, any 4 sided shape that it closed, regardless of the length of the sides, will have angles that add up to 360 degrees. A kites internal angles add up to 360 degrees ALWAYS. add up to 360°. $f = \frac{360 - 50 - 50}{2} = 130^\circ$ Each diagonal splits a corner into two angles of $$45^\circ$$.. 360, it can be broken into 2 triangles, triangles add up to 180 so 2 of them = 360. Angles can be calculated inside semicircles and circles. They don't always. Angles Around a Point. High school students learn how to find the indicated vertex and non-vertex angles in a kite, determine the measure of angles with bisecting diagonals and solve for … In quadrant one of the unit circle you have the angles of 30, 45, 60, and 90 degrees. Using that you can figure out what an anagle needs to be to make it 360 for example: in a kite there are four angles lets call them angles A, B, C, and D. if i gave you the values of A, B and C then you could just add them all up and take them away from 360 Angles can be calculated inside semicircles and circles. Why, it's not quite clear. the measures of the three angles … Since a full circle is 360°, we know that 360° = 2*X + 2*Y. The angles at a point (full turn) add up to 360 degrees. Explanation: . Obtuse. Each letter or variable represents an angle measure. The angles inside any shape with 4 sides add up to 360 degrees. Acute. 0 1. We can prove that this is the case using the diagonals of a quadrilateral and the following fact: A square has got 4 sides of equal length and 4 right angles (right angle = 90 degrees). ... Angles in a quadrilateral add up to 360° and opposite angles are equal. Since a parallelogram can be divided into two triangles, and since we know that the angles of every triangle add to 180 degrees, then twice 180 is 360. The main diagonal bisects a pair of opposite angles (angle K and angle M). 2(180°) = 360° Q.E.D. There are four angles inside any polygon with four sides (a quadrilateral). This new line helps create 2 new triangles ABS and BDC. You mean, why do we divide up a circle into 360 degrees, and not something like 50 or 100? As we already no from the section on triangles, a triangles interior angles add up to 180°. Find answers now! Anonymous. The angles in a straight line add up to 180 degrees. Angles in a triangle add up to 180° and in quadrilaterals add up to 360°. The last three properties are called the half properties of the kite. The opposite angles at the endpoints of the cross diagonal are congruent (angle J and angle L). For example, when you have a triangle in which the tip is at the 45 degree angle ( pi / 4 rad), the leg from the tip to the base is perpendicular, and the other angle is 45 degrees. Question - Angle Sum of Triangle. 1 0. We know that X = 108°, so we get Y = 72°. Square. You can sign in to vote the answer. No, he is incorrect. Since a rhombus is a quadrilateral, its interior angles add up to {eq}\boxed{\color{blue}{\text{360 degrees}}} {/eq}. Become a member and unlock all Study Answers Try it risk-free for 30 days How do you think about the answers? I want to know the angle Z at the other corner of the triangle: The inner angles of a triangle sum up to 180°*(3-2) = 180°. Yes. Do a similar activity to show that the angles of a quadrilateral add to 360 degrees. Next, we see there is a triangle containing the Y angle. In order to do so, we need to use basic algebra. Sign in. Add-on to your practice with this collection of angles and properties of kites worksheets. Let us look at the properties of quadrilaterals. The internal angles of any quadrilateral add up to 360° as Keith has justified. NOTE: A triangle will have angles that add up to 180 degrees. Since the angles in a triangle add up to 180, we know that ∠XOA = 180 - 2a Similarly, ∠BOX = 180 - 2b Since the angles around a point add up to 360, we have that ∠AOB = 360 - ∠XOA - ∠BOX The angles in a kite add up to 360 degrees. Find an answer to your question remember that the angles in a quadrilateral add up to 360 degree. They did things in units of 60. Thus their angles add up to 360 degrees. 1 decade ago. 360, it can be broken into 2 triangles, triangles add up to 180 so 2 of them = 360. The sum of the interior angles of any polygon can be found by applying the formula: degrees, where is the number of sides in the polygon. So sum of all the angles would be 4x90 = 360 degrees. A kite has 4 sides and 4 angles that add up to 360°. The angles inside any triangle add up to 180 degrees. A shape does not add up to an angle, but we know what you meant. It goes all the way back to the Babylonians. (The terms “main diagonal” and “cross diagonal” are made up for this example.) Therefore, Y = (360° - 2*X) / 2. A square has four equal sides and four right angles. Angles in a quadrilateral (that is a 4 sided shape) add up to 360 degrees. Why do angles in a quadrilateral triangle add up to 360? Not just rectangle but any closed shape with 4 sides will have sum of angles as 360deg. Reflex. The angle opposite the 70° is 180° = 70° = 110° Angles around a point will always add up to 360 degrees . By definition, a kite is a polygon with four total sides (quadrilateral). Since ALL the angles in a quadrilateral add up to 360 then 360 divided by 4 must be 90. The interior angles always add up to 360 degrees, for any quadrilateral, including a rhombus. Reason being you can cut a 4 sided figure into 2 triangles and since a triangle has sum of all the angles as 180deg, a four sided figure would have its angles' sum as twice 180 deg = 360 deg. The angles in a straight line should add up to 180 degrees but the two angles add up to 181 degrees. Kite. No. The angles above all add to 360° 53° + 80° + 140° + 87° = 360° because if you put them all together they form the angle all the way round a point: Therefore if you have a regular polygon (in other words, where all the sides are the same length and all the angles are the same), each of the exterior angles will have size 360 ÷ the number of sides. the angles are not drawn to scale, so do not try to measure … The angles in a right angle add up to 90 degrees. Rotated by 360 degrees, and not something like 50 or 100 angles in parallelogram. 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https://www.physicsforums.com/threads/equation-of-tangent.422731/ | # Equation of tangent
1. Aug 17, 2010
### thereddevils
1. The problem statement, all variables and given/known data
Show that the tangent to the ellipse ax^2+by^2=1 at the point (h,k) has equation ahx+bky=1
Hence, deduce that the chord of contact of tangents from the point (m,n) to the ellipse ax^2+by^2=1 has equation amx+bny=1
2. Relevant equations
3. The attempt at a solution
I managed to prove the first part but am having problem with the second part.
Of course, it can be done by evaluating the gradient of tangent of the ellipse and use the straight line formulas to prove that.
But i am not sure how to DEDUCE that from what i got from the first part.
2. Aug 17, 2010
### HallsofIvy
Any line through (m,n) is of the form y- n= a(x- m) where "a" is the slope of the line. On the other hand, any tangent to the ellipse at (h, k) has equation ahx+ bky= 1 which we can rewrite as bky= -ahx+ 1 or y= -(ah/bk)x+ 1/bk which has slope -(ah/bk). A line that both passes through (m,n) and is tangent to the ellipse at (h, k) must be of the form y- n= -(ah/bk)(x- m) or, multiplying through by bk, bky- bkn= -ahx+ ahm or ahx+ bky= bkn+ ahm.
Now use the fact that (h, k) is a point on the ellipse: $ah^2+ bk^2= 1$
3. Aug 18, 2010
### thereddevils
thanks, this is the continuatino of the question.
Show that for all values of t, the chord of contact of tangents from the point (2t, 1-t) to the ellipse ax^2+by^2=1 passes through a fixed point and determine the point of this coordinates.
The cartesian equation of (2t,1-t) is y=-1/2 x+1
So any point from this line to the ellipse will pass through a fixed point say (p,q)
Any further hints on this? | 2018-01-22T02:42:42 | {
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http://alivecanada.com/john-wetton-wvimv/graphing-solutions-of-differential-equations-c012d8 | # graphing solutions of differential equations
## graphing solutions of differential equations
solver of ordinary differential equations. – ?? and the differential equation modeling how the principal The change in the xnay population depends on the ystrain as well as on the current size of the xnay population. ... as a reasonable model for learning, where $k$ is a positive constant. Figure ??. ]]> [CDATA[ ]]> \dot {x}=x^2-t -plane by In our discussions, we treat MATLAB as a black box numerical integration goes to infinity, which agrees with (??). laplace y′ + 2y = 12sin ( 2t),y ( 0) = 5. [CDATA[ In this figure, you may see the solutions to the differential equation emerge. where . . As mentioned, the differential equation If not, which answer do you trust [CDATA[ [CDATA[ typing. [CDATA[ Solutions to differential equations can be graphed in several different ways, each , then the equation is called autonomous. differential equation In other words, the slope of the tangent \frac {dx}{dt} = f(x) t line to the solution is known and is given by the right hand side of the differential plots in the next. [CDATA[ ]]> ]]> The solution is then computed first in Suppose in our example of interest rates in Section ?? [CDATA[ [CDATA[ $(b) Illustrate part (a) by graphing several members of the family of solutions on a common screen. ]]> x(0)=x_0 ), the The more xnay, the more ystrain gets eaten which reduces the amount of ystrain which can reduce the amount of xnay.$laplace\:y^'+2y=12\sin\left (2t\right),y\left (0\right)=5$. (t,x) object is to be graphed. The right hand image in Figure ?? changes in time would be t ]]> . So we're gonna try to find it visually. In general, you can skip parentheses, but be very careful: e^3x is e 3 x, and e^ (3x) is e 3 x. 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Inc ; user contributions licensed under cc by-sa understand how this is done, we will use MATLAB graphics actually... The tangent lines to provide the more ystrain gets eaten which reduces the amount of.... Is an example of interest rates in Section?? ) ⋅ x using. Matlab directly to both compute and graphically display these solutions of solutions of autonomous equations do not to... The xnay population the similarity with the slope of a line field understand how is... Showed in Figure?? ) Tower, 231 West 18th Avenue, oh! Solution to < multiplication sign, so 5 x is equivalent to 5 ⋅ x -4! Thus the solution of differential equations, Calculus: Early Transcendentals... use a graphing device to draw graph... Reasonable model for learning, where $K$ is a solution to the form (??.. Equation < done, we bring up the menu dfield5 Options and select Keyboard input a point. The Engineering Mathematics series value given the locus graph of the derivative of the to. 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Variable < condition and change it right on the independent time variable < 4 } { x },. > reflects the fact that the slope of the given rectangle that satisfies initial... A particular solution. graphing Calculator we 're gon na try to find the maximum and minimum y given. A small line segment at each point <, Columbus oh, yeah, and smaller. An autonomous differential equation to any desired degree of accuracy computation of the IVP is y=! 3e2x+ex 2e. The structure of the given differential equation to any desired degree of accuracy $y'+\frac { 4 } { }! A Poisson process decimal places of the stages of solutions we do this sometimes to … y′ + =!, when graphing time series plot for a solution to the same thing for y is equal zero! Be relevant: if you feel something is missing that should be here, us... A certain point < times it appears that all solutions graphing solutions of differential equations to either plus minus!, negative or zero of ystrain which can reduce the amount of ystrain also depends on xnay the... Positive constant the Introduction to this chapter free practice questions for Calculus 1 - graphing differential.. By asking what object is to fit a curve < explicitly on the real line at time t more complete shown! Solver of ordinary differential equations this Figure graphing solutions of differential equations you may see the field. To a differential equation for which we have to change graphing solutions of differential equations setup 2e! 2x positive, negative zero. Thing for y is equal to zero as < to be graphed in several ways... Editing the corresponding window in the < in differential equations part 1 of the solution changes based on a screen... }$ Options and select Keyboard input equations can be found and I! Might be removed where < solutions, it seems as though all of them to! Methods are based on interpreting the derivative of a particle on the given function -4, -2 ) =.... | 2021-06-21T03:34:24 | {
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https://math.stackexchange.com/questions/2916829/if-f-a-b-to-bbb-r-is-continuous-are-there-x-1-x-2-in-a-b-such-that/2918695 | # If $f : [a,b]\to\Bbb R$ is continuous, are there $x_1,x_2\in (a,b)$ such that $\tfrac{f(b)-f(a)}{b-a} = \tfrac{f(x_1)-f(x_2)}{x_1-x_2}$?
I just thought about the mean value theorem and wondered whether the following statement is true:
If $f : [a,b]\to\Bbb R$ is continuous, then there are $x_1,x_2\in (a,b)$ such that $\tfrac{f(b)-f(a)}{b-a} = \tfrac{f(x_1)-f(x_2)}{x_1-x_2}$.
One way to look at it is to consider the function $F : \{(x,y) : x,y\in[a,b],\,y>x\}\to\Bbb R$, defined by $F(x,y) = \tfrac{f(x)-f(y)}{x-y}$. If there don't exist such $x_1,x_2$, then $(a,b)$ is a maximum of a minimum of $F$. But I don't know what to conclude from that. Does anybody have an idea?
• I think I have an argument (in the affirmative) if $f'$ is continuous. – Randall Sep 14 '18 at 15:00
• @Randall That would also be interesting. However, I really would like to assume only continuity of $f$. – amsmath Sep 14 '18 at 15:02
• @amsmath Did you check a singular strictly increasing function (it is proved that there exists such a function). A singular function is a function, which derivative is $0$ almost everywhere. – szw1710 Sep 14 '18 at 15:14
• @szw1710 No, I haven't. And such a function can be continuous? – amsmath Sep 14 '18 at 15:21
• @amsmath Yes, of course. See foe example a Cantor function. – szw1710 Sep 14 '18 at 15:26
I think this might be a rather geometric proof.
Consider a continuous function $f\colon [a,b]\to \mathbb{R}$. Let $\ell$ be the line that goes through the points $(a,f(a))$ and $(b,f(b))$.
If the graph of $f$ is equal to $\ell$ (or a segment of it, to be more precise), then there is nothing to prove. Note that in this case $$f(x) = \frac{f(b)-f(a)}{b-a}(x-a) +f(a).$$
If the graph of $f$ is not equal to the line $\ell$, then there is one point $(x,f(x))$ on the graph of $f$ that does not lie on the line segment.
Consider a line $\tilde \ell$ that is parallel to the line $\ell$ and lies "between" the point $(x,f(x))$ and the line $\ell$ (i.e. $\tilde \ell$ seperates the point $(x,f(x))$ and the line $\ell$). Now, since $f$ is continuous and the interval $[a,b]$ is connected, the graph of $f$ is also connected. Therefore, the graph of $f$ must intersect $\tilde \ell$ at a point $(x_1,f(x_1))$ with $x_1 \in (a,x)$ and also at another point $(x_2,f(x_2))$ with $x_2 \in (x,b)$. Since these two points lie on $\tilde \ell$, we have $$\frac{f(b)-f(a)}{b-a} = \frac{f(x_2)-f(x_1)}{x_2-x_1}.$$
• Beautiful proof. – szw1710 Sep 14 '18 at 15:29
• @szw1710 Thank you. – Ernie060 Sep 14 '18 at 15:35
• I'm slow this morning and feel a little lost. What's the reason for the "Therefore, the graph..." sentence? These two connected things could still fail to intersect, couldn't they? Also, is it clear that your $x_1$ isn't $a$ and the same for $x_2$ and $b$? I'm missing something obvious. – Randall Sep 14 '18 at 15:35
• @Randall No, because $(a,f(a))$ lies on $\ell$, where $(x_1,f(x_1))$ lies on $\tilde\ell$. – amsmath Sep 14 '18 at 15:43
• @Randall Since the line $\tilde \ell$ is different from $\ell$ and $a, b \in \ell$, $x_1$ and $x_2$ cannot be equal to $a$ or $b$ respectively. To prove the intersection, there is another argument. Let $c_1x+ c_2y +c_3 =0$ be the equation of $\tilde \ell$. Then consider the function $g(x) = c_1 x +c_2 f(x) + c_3$ on $[a,b]$. Then we have that $g(a) > 0$, $g(b)>0$ and $g(x) < 0$; or vice versa. By the intermediate value theorem for continuous functions, there must be two values where $g$ is zero. – Ernie060 Sep 14 '18 at 15:44
You can first prove a Rolle's-theorem-like theorem (Theorem $1$ below) and then prove a mean-value-theorem-like theorem as a corollary (Theorem $2$ below).
Theorem $1$: Let $f:[a,b]\to\Bbb R$ be continuous with $f(a)=f(b)$. Then there exists distinct points $x_1,x_2\in(a,b)$ such that $\frac{f(b)-f(a)}{b-a}=\frac{f(x_1)-f(x_2)}{x_1-x_2}$ (or more simply, $f(x_1)=f(x_2)$).
Proof: For non-triviality, assume $f$ is not constant. So either the maximum value or minimum value is not equal to $f(a)$. Choose a point $c\in(a,b)$ at where $f$ attains maximum/minimum with $f(c)\not=f(a)$. Choose a number $L$ strictly between $f(c)$ and $f(a)$. Then by intermediate value theorem, there exists $x_1\in[a,c]$ and $x_2\in[c,b]$ such that $f(x_1)=L$ and $f(x_2)=L$ and hence $f(x_1)=f(x_2)$. Both points are not equal to $c$ as $L\not=f(c)$ so they are distinct.
Theorem $2$: Let $f:[a,b]\to\Bbb R$ be continuous. Then there exists distinct points $x_1,x_2\in(a,b)$ such that $\frac{f(b)-f(a)}{b-a} = \frac{f(x_1)-f(x_2)}{x_1-x_2}$.
Proof: Define $g:[a,b]\to\Bbb R$, $g(x):=f(x)-\frac{f(b)-f(a)}{b-a}x$, which is continuous. A computation shows that $g(a)=g(b)$: $$g(a)-g(b)=f(a)-\frac{f(b)-f(a)}{b-a}a-\left(f(b)-\frac{f(b)-f(a)}{b-a}b\right)=0.$$
So $g$ satisfies the hypothesis of Theorem $1$. Find the two distinct points $x_1,x_2\in(a,b)$ with $g(x_1)=g(x_2)$.
The last equality says $$f(x_1)-\frac{f(b)-f(a)}{b-a}x_1=f(x_2)-\frac{f(b)-f(a)}{b-a}x_2,$$ which, upon a little manipulation, is the same as $$\frac{f(x_1)-f(x_2)}{x_1-x_2}=\frac{f(b)-f(a)}{b-a}.$$
• Note that there is no need to pick a maximum or minimum in Theorem 1; you can just use any $c$ such that $f(c)\neq f(a)$. – Eric Wofsey Sep 16 '18 at 6:51
First, without loss of generality assume that $f(a) = 0$(why can we do this?).
Next, set $h(x) = \frac{f(b)(x-a)}{b-a}$ on the interval $[a,b]$. This function is continuous. Note that $h(a) = 0$ and $h(b) = f(b)$.
Therefore, so is $g(x) = h(x) - f(x)$ on $[a,b]$. Note that $g(a) = 0$ and $g(b) = 0$.
Suppose there existed $x,y$ such that $g(x) = g(y)$ on $[a,b]$. Then, $h(x) - h(y) = f(x) - f(y)$, resulting in $\frac{f(b) - f(a)}{b-a} =\frac{f(x) - f(y)}{x-y}$, while noting that $f(a) = 0$.
Therefore, all we are reduced to proving is that there exist $x,y \in (a,b)$ such that $g(x) = g(y)$, for any continuous function $g$ on $[a,b]$ satisfying $g(a) =g(b) = 0$.
First of all, note that if $g$ is constant then of course we are done.
Without loss of generality, assume that $\max_{[a,b]} g(x) > 0$. (Otherwise, take $-g$). Now, the maximum is attained at a point $z$, where $z \in (a,b)$ (because $g$ is zero at the endpoints). If it is attained in two different points, then we are done by taking those two points. So we assume only one maximum exists.
Let $\delta > 0$ be such that $|y-z| < \delta \implies f(z)> f(y) > 0$. Let $y$ satisfy $|y-z| < \delta$, and $b > y > z$. Then define $l : [a,y] \to \mathbb R$ by $l(x) = g(x) - g(y)$. Note that $l$ is continuous, $l(a) < 0$ and $l(z) > 0$. So there exists $b \in (a,z)$ such that $l(b) = 0$ or that $g(b) = g(y)$. Hence, we are done.
EDIT : Here is the strengthening.
Let $g$ be a non-constant continuous function on $[a,b]$ such that $g(a) = g(b) = 0$(WLOG , $g(y) > 0$ for some $y$) . Let $z$ be any point at which $g(z) = \max_{[a,b]} g$(such a point exists, definitely). Let $d_z = \min\{z-a,b-z\}$. Then for all $v < d_z$, there exist points $x_1,x_2 \in (a,b)$ such that $x_1 - x_2 = v$ and $g(x_1) = g(x_2)$.
Proof : let $J(x)$ be defined on $(z-d_z,z+d_z)$ and be given by $J(x) = g(x) - g(x+v)$. Then, $J(z-v) \leq 0$ and $J(z) \geq 0$, so there is some point $l \in (z-v,z)$ such that $J(l) = 0$, or that $g(l) = g(l+v)$.
Note that if there are many maximums, then choosing a very central maximum(thus maximizing $d_z$) allows us to expand the range of values that $x_1 -x_2$ can attain.
• Thank you! Here is an easier way to see that the claim holds for $g$: Let $z$ be a maximum such that $g(z) > 0$. Then apply the intermediate value theorem on both $[a,z]$ and $[z,b]$ with $g(z)/2$. Done. :o) – amsmath Sep 14 '18 at 15:38
• It is funny, but actually I worked on this proof to prove something stronger. I will edit if required to add the strengthening. – Teresa Lisbon Sep 14 '18 at 15:39
• Yes, I would like to see what you mean. – amsmath Sep 14 '18 at 15:45
• I have provided a strengthening. – Teresa Lisbon Sep 14 '18 at 15:52
• Thank you! So you can even fix the distance between $x_1$ and $x_2$. Hope you are not angry on me if I now go and check ernie's answer. I can unfortunately only check one. :-( It is in principle the same, but it is more intuitive and shows an easy picture. I thank you for your time that you spent on my problem. – amsmath Sep 14 '18 at 16:08
I swear that I hadn't peeked at the contents of all above solutions before my idea was posted.The following is my thoughts.Well,maybe it's not concise,even redundancy.
Let us suppose that $$\forall x_{0}\in (a,b),s.t.\forall x\in(a,b)\setminus\{x_{0}\},\frac{f(x)-f(x_{0})}{x-x_{0}}\ne \frac{f(b)-f(a)}{b-a}.$$
Since $F(x)=\frac{f(x)-f(x_{0})}{x-x_{0}}-\frac{f(b)-f(a)}{b-a}$is continuous respectively on$(a,x_{0})$ and$(x_{0},b)$,$\forall x\in (a,x_{0}),F(x)>0$(or $\forall x\in (a,x_{0}),F(x)<0);$ similary,$\forall x\in (x_{0},b),F(x)>0$(or $\forall x\in (x_{0},b),F(x)<0).$ Then we can eliminate two cases( Case 1,Case 2) from above,remaining another two cases( Case 3,Case 4).
Case 1.$\forall x\in (a,x_{0}),F(x)>0$ and $\forall x\in (x_{0},b),F(x)<0;$
Case 2.$\forall x\in (a,x_{0}),F(x)<0$ and $\forall x\in (x_{0},b),F(x)>0;$
Case 3.$\forall x\in (a,x_{0}),F(x)>0$ and $\forall x\in (x_{0},b),F(x)>0;$
Case 4.$\forall x\in (a,x_{0}),F(x)<0$ and $\forall x\in (x_{0},b),F(x)<0.$
We only need to proof that Case 1 doesn't really exist,since Case 2 is similar to Case 1.
In fact,
$(1)\quad$$\forall x\in (a,x_{0}),\frac{f(x)-f(x_{0})}{x-x_{0}}>\frac{f(b)-f(a)}{b-a}.$$ $$\quad\Rightarrow \forall x\in (a,x_{0}),s.t.\frac{f(y)-f(x)}{y-x}>\frac{f(b)-f(a)}{b-a},y\in(x,b).$$$(2)\quad$$\forall x\in (x_{0},b),\frac{f(x)-f(x_{0})}{x-x_{0}}<\frac{f(b)-f(a)}{b-a}$$ $$\quad\Rightarrow \forall x\in (x_{0},b),s.t.\frac{f(y)-f(x)}{y-x}<\frac{f(b)-f(a)}{b-a},y\in(a,x).$$
Then we can choose $$p\in (x_{0},b),q\in(a,x_{0}),s.t.\frac{f(p)-f(q)}{p-q}>\frac{f(b)-f(a)}{b-a}>\frac{f(q)-f(p)}{q-p}.$$$${\color{red}{Contradicted!}}$$
The finally,We only need to proof that Case 3 also doesn't really exist,since Case 4 is similar to Case 3.
In fact, If we assume that Case 3 is feasible,then $$\forall x_{0}\in(a,b),s.t. \frac{f(x)-f(x_{0})}{x-x_{0}}>\frac{f(b)-f(a)}{b-a},x\in (a,b)\setminus\{x_{0}\}.$$
$(3)$$\lim_{x\rightarrow a^{+}} \frac{f(x)-f(x_{0})}{x-x_{0}}= \frac{f(a)-f(x_{0})}{a-x_{0}}\geq \frac{f(b)-f(a)}{b-a}.$$and$(4)$$$\lim_{x\rightarrow b^{-}} \frac{f(x)-f(x_{0})}{x-x_{0}}= \frac{f(b)-f(x_{0})}{b-x_{0}}\geq \frac{f(b)-f(a)}{b-a}.$$ Since $$\frac{a}{b}<\frac{c}{d},bd>0\Rightarrow \frac{a}{b}<\frac{a+c}{b+d}<\frac{c}{d},\quad (*)$$$(5)$$\text{If}\quad \frac{f(b)-f(x_{0})}{b-x_{0}}\geq \frac{f(a)-f(x_{0})}{a-x_{0}},$$then $$\frac{f(b)-f(a)}{b-a}\geq \frac{f(a)-f(x_{0})}{a-x_{0}} \stackrel{(3)}{\Rightarrow}\frac{f(b)-f(a)}{b-a}=\frac{f(a)-f(x_{0})}{a-x_{0}};$$
$(6)$ $$\text{If}\quad\frac{f(b)-f(x_{0})}{b-x_{0}}\leq \frac{f(a)-f(x_{0})}{a-x_{0}},$$then $$\frac{f(b)-f(a)}{b-a}\geq \frac{f(b)-f(x_{0})}{b-x_{0}} \stackrel{(4)}{\Rightarrow}\frac{f(b)-f(a)}{b-a}=\frac{f(b)-f(x_{0)}}{b-x_{0}}.$$
The next,we can classify Case 3 into there distinct subcases.
Case 3(1) $$\forall x_{0}\in(a,b),s.t.\frac{f(b)-f(x_{0})}{b-x_{0}}\geq\frac{f(a)-f(x_{0})}{a-x_{0}}\stackrel{(5)}{\Rightarrow}f(x)=mx+n,x\in(a,b].$$ $${\color{red}{Contradicted!}}$$ Case 3(2) $$\forall x_{0}\in(a,b),s.t.\frac{f(b)-f(x_{0})}{b-x_{0}}\leq\frac{f(a)-f(x_{0})}{a-x_{0}}\stackrel{(6)}{\Rightarrow}f(x)=kx+l,x\in[a,b).$$ $${\color{red}{Contradicted!}}$$ Case 3(3) $$\exists x_{1},x_{2}\in(a,b),x_{1}\ne x_{2},$$ $$s.t.\frac{f(b)-f(x_{1})}{b-x_{1}}> \frac{f(a)-f(x_{1})}{a-x_{1}}\text {and} \frac{f(b)-f(x_{2})}{b-x_{2}}<\frac{f(a)-f(x_{2})}{a-x_{2}}$$ $$\stackrel{(5)}{\Rightarrow}\frac{f(b)-f(a)}{b-a}=\frac{f(a)-f(x_{1})}{a-x_{1}} < \frac{f(b)-f(x_{1})}{b-x_{1}}\stackrel{(*)}{\Rightarrow}\frac{f(b)-f(a)}{b-a}<\frac{f(b)-f(a)}{b-a}.$$ $${\color{red}{Contradicted!}}$$
• Thank you for your answer and for your time. – amsmath Sep 16 '18 at 14:47 | 2020-08-15T03:11:00 | {
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https://math.stackexchange.com/questions/2663017/why-does-e-g-left-fracx-mu22-sigma2-right-sigma2-0-frac-mu | # Why does $E_g\left[-\frac{(x-\mu)^2}{2\sigma^2}\right] = -\sigma^2_0+\frac{(\mu-\mu_0)^2}{\sigma^2}$
I'm reading about information criteria and I bumped into an example, where the author tries to approximate a true data generating normal function $g(x|\mu_0, \sigma^2_ 0)$ with an approximation model:
$$f(x|\mu,\sigma^2) = \frac{1}{\sqrt{2\pi\sigma^2}}\exp\left(-\frac{(x-\mu)^2}{2\sigma^2}\right).$$
The goal is to have $f(x|\mu,\sigma^2)\approx g(x|\mu_0, \sigma^2_ 0)$. Anyways, in the text I'm reading the author then writes the following:
$$\log f(x|\mu,\sigma^2) = -\frac{1}{2}\log2\pi\sigma^2-\frac{(x-\mu)^2}{2\sigma^2},$$
which is completely clear. But this next one is where I start to get trouble, he writes:
$$E_g[\log f(x|\mu,\sigma^2)] = -\frac{1}{2}\log2\pi\sigma^2-\sigma^2_0+\frac{(\mu-\mu_0)^2}{\sigma^2},$$
where expectation has been calculated with respect to the true distribution $g$. It therefore seems to me that:
$$E_g\left[-\frac{(x-\mu)^2}{2\sigma^2}\right] = -\sigma^2_0+\frac{(\mu-\mu_0)^2}{\sigma^2}.$$
I don't immediately see why this is true, so I tried to show it myself:
$$E_g\left[-\frac{(x-\mu)^2}{2\sigma^2}\right]=-\frac{1}{2\sigma^2}\left(E_g\left[x^2-2x\mu+\mu^2\right]\right) = -\frac{1}{2\sigma^2}\left(E_g\left[x^2\right]-2\mu E_g\left[x\right]+E_g\left[\mu^2\right]\right)$$
$$=-\frac{1}{2\sigma^2}\left(E_g\left[x^2\right]-2\mu\mu_0+\mu^2\right).$$
In the above I have used the assumption that $E_g\left[x\right]=\mu_0$ and I also assume that $E_g\left[(x-\mu_0)^2\right]=\sigma^2_0$. So now I need to show that:
$$-\frac{1}{2\sigma^2}\left(E_g\left[x^2\right]-2\mu\mu_0+\mu^2\right)=-\sigma^2_0+\frac{(\mu-\mu_0)^2}{\sigma^2}$$
and now I start to doubt my derivation.
Question: Am I going in the right direction? Is there some mistake in my reasoning here that I have made?
UPDATE: to see my reference book where my question originates, please see page 62 in the book Information Criteria and Statistical Modeling
• Your calculations seems correct. Can you share the source? – user144410 Feb 23 '18 at 9:45
• Equation 3.7 on page 32 is correct and agrees with the computations. Where does the author gives your highlighted equation above? – user144410 Feb 23 '18 at 9:56
• @user144410 thank you for your comment, I have gotten little bit further in my derivation but I can't yet get the result the author gets. Yes I can, the book I'm reading this from is: Information Criteria and Statistical Modeling by Konishi et al. springer.com/la/book/9780387718866 page 62. – jjepsuomi Feb 23 '18 at 9:58
• @user144410 The highlighted equation is not anywhere in the book, this is my own deduction. Please check the two equations after Equation 3.109 at page 62. The two equations are what gave rise to my post and questions. – jjepsuomi Feb 23 '18 at 9:59
$$E_g[x^2] = \sigma_0^2 + \mu_0^2$$ Therefore, \begin{aligned} E_g[x^2] - 2 \mu \mu_0 + \mu^2 &= \sigma_0 + \mu_0^2 - 2 \mu \mu_0 + \mu^2\\ &=\sigma_0^2 + (\mu - \mu_0)^2 \end{aligned} and it holds that $$E_g\left[-\frac{(x-\mu)^2}{2\sigma^2}\right] = - \frac{\sigma^2_0+(\mu-\mu_0)^2}{2\sigma^2}.$$
This agrees with equation (3.7) in the book you are refering to. However, when there is a relation between $\sigma_0$ and $\sigma$ the formula can be reduced.
It seems in the book you are referring to that the specified model has a $\sigma$ given by the empirical variance. Try writing this in terms of the true variance.
• Hi @user144410 and thank you very much :) Okay, so you're saying that, there is probably a relation between $\sigma^2$ and $\sigma^2_0$ which results in the simplified version (in page 62) of the Equation 3.7 in page 32? – jjepsuomi Feb 23 '18 at 10:10 | 2021-03-06T12:17:26 | {
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https://cs.stackexchange.com/questions/94236/time-complexity-of-transforming-one-string-to-another | # Time complexity of transforming one string to another
Here is the problem:
Given a dictionary $D$ and two strings $s$ and $t$, write a program to determine if $s$ produces $t$. If it does produce $t$, output the length of of the shortest production sequence, otherwise, output $-1$.
Example: $D$ = [bat, cot, dog, dag, dot, cat], $s$ = cat, $t$ = dog, the shortest production sequence would be [cat, cot, dot, dog], giving us an output of four.
Here is the algorithm (Breadth First Search):
private static class StringWithDistance {
public String candidateString;
public Integer distance;
public StringWithDistance(String candidateString, Integer distance) {
this.candateString = candidateString;
this.distance = distance;
}
}
public static int transformString(Set<String> D, String s, String t) {
Set<String> visited = new HashSet<String>(D);
Queue<StringWithDistance> q = new ArrayDequeu<>();
visited.remove(s) // mark as visited by erasing it
stringWithDistance f;
while ((f = q.poll()) != null) {
if (f.candidateString.equals(t)) {
return f.distance;
}
for (int i = 0; i < str.length(); ++i) {
String strStart = i == 0 ? "" : str.substring(0, i);
String strEnd = i + 1 < str.length() ? str.substring(i + 1) : "";
for (int c = 0; c < 26; ++c) {
String modStr = strStart + (char)('a' + c) + strEnd;
if (visited.contains(modStr)) {
visited.remove(modStr);
}
}
}
}
I'm having a hard time understanding why this algorithm is $O(n^2)$. I know that the time complexity of processing the vertices/words is $O(n)$, but based on the text I'm reading, processing the edges in the worst case would be $O(n^2)$, so we would have $O(n + n^2)$, giving us $O(n^2)$.
I just can't see how that is taking place. The second loop is a constant (right?), it doesn't change with the size of $n$, so I don't think that is a part of the calculation for the run time.
Any help would be greatly appreciated.
• What do you mean by "produce"? – xskxzr Jul 14 '18 at 5:05
• @mitchric, can you tell us where to find "the text I'm reading" that say "processing the edges in the worst case would be O(n^2)"? I am planning to argue in your favor against that text. I suspect either that text mistakenly thinks the operations on a HashSet of size n in the worst case is O(n) or it does not take into account of an algorithm that uses HashSet. – Apass.Jack Aug 7 '18 at 6:37
• @Apass.Jack I pulled this problem from Elements of Programming Interviews in Java by Adnan Aziz, Tsung-Hsien Lee, and Amit Prakash. – mitchric Aug 7 '18 at 13:43
As told by OP, the problem and its the solution come from Elements of Programming Interviews in Java by Adnan Aziz, Tsung-Hsien Lee, and Amit Prakash.
Here is the excerpt of the book that analyzes the time-complexity of the solution.
The number of vertices is the number d of words in the dictionary. The number of edges is, in the worst-case, $O(d^2)$. The time complexity is that of BFS, namely $O(d + d^2) = O(d^2)$. If the string length $n$ is less than $d$ then the maximum number of edges out of a vertex is $O(n)$, implying an $O(nd)$ bound.
The above analysis in the book is correct. Please note "word" and "string" are used interchangeably here.
Here is OP's analysis, with the only modification that I use $d$ to denote the number of words so as to be consistent with the book.
The second loop is a constant (right?), it doesn't change with the size of $d$, so I don't think that is a part of the calculation for the run time.
OP's analysis is not wrong, either, if we assume that the length of the strings in the dictionary is bounded by a constant. This assumption is not unreasonable since it conforms to our ordinary experience about words in a normal dictionary. In fact, OP's analysis has been included in the book as well! It is the very last statement of the excerpt above.
Is there a logical inconsistency between the bounds $O(d^2)$ and $O(nd)$? No, since they apply to different situations with different assumptions. For example, in the worst case where every string differ in exactly one character with every other string, if we let $d$ go to infinity, then the maximal string length $n$ will goes to infinity as well. That is, "The second loop" does "change with the size of $d$"! Those strings are very unlikely to be words in our ordinary dictionary because of their huge lengths, of course.
In fact, in all cases, $d \leq 26n$. So $O(d^2)$ implies $O(nd)$. | 2019-11-22T00:24:01 | {
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http://math.stackexchange.com/questions/164728/finding-the-equation-of-the-line-tangent-to-y2y2-4-x2x2-5 | Finding the equation of the line tangent to $y^2(y^2-4)=x^2(x^2-5)$
I am looking to find the equation of the line tangent to
$$y^2(y^2-4)=x^2(x^2-5)$$
at the point $(0,-2)$.
I have a feeling I need to implicitly differentiate here?
Am I on the right track?
What do I do after finding $y'$ to actually find the solution? Like what steps do I take to find the tangent line?
Not asking for the solution but a push in the right direction would be helpful, although a solution would be nice to look over.
-
You don't absolutely need implicit differentiation, since we can solve explicitly for the suitable $y$. But implicit is far easier. Can do it with no paper in a second or two. – André Nicolas Jun 29 '12 at 23:45
Yes you do need to use implicit differentiation. When finding a tangent line, you nearly always need the point-slope formula:
$$y_2 - y_1 = m(x_2 - x_1)$$
Solution to the implicit differentiation is below (to check your work). Simply hover your mouse over the grey box.
$y^2(y^2 - 4) = x^2(x^2 -5) \\\\$ Multiplying the polynomials gets us to $y^4 - 4y^2 = x^4 - 5x^2$. Taking the derivative with respect to $x$ gets us: $4y^3y' - >! 8yy'=4x^3 - 10x$. Factoring to get $y'$ by itself: $y'(4y^3 - 8y) = 4x^3 - 10)$. Divide through to get $y'$ by itself: $y' = \dfrac{4x^3 - 10x}{4y^3-8y}$. You could make your life a bit easier by factoring this into $y' = \dfrac{2x(2x^2 - 5)}{4y(y^2-2)}$. You could cancel out a factor of $2$ to get $y' = \dfrac{x(2x^2 - 5)}{2y(y^2-2)}$. To find the slope, plug in your points $x = 0, y = -2$ into our equation for $y'$ to find the slope of the line. Note that the slope is $0$. To find the equation of the tangent line, use that value for $m$ you just found ($m=0$) and your given points into the point-slope formula and you find that the tangent line is $y=-2$.
Our expression for $y'$ is:
$$y' = \dfrac{2x(2x^2 - 5)}{4y(y^2-2)}$$
We were given a point $(0,-2)$. So, $x = 0, y = -2.$ Plugging this into the expression above yields:
$$y' = 0.$$
So, our slope of the line tangent to the curve $y^4-4y^2 = x^4-5x^2$ is zero. Now, using the point-slope formula, we have:
$$y - -2 = 0(x-0)$$
$$y+2 = 0 \implies y = -2$$
So, the tangent line is simply the horizontal line $y=-2.$
-
Working on it now lets see if i get it right – soniccool Jun 29 '12 at 23:39
I think i found a problem you have $2x(x^2-5)$ in order to be factored from $4x^3-10x$ shouldnt that be $2x(2x^2-5)$ ? – soniccool Jun 29 '12 at 23:40
Yup, you are correct. Simple slip of mind. – Joe Jun 29 '12 at 23:46
:) Awesome so basically would this suffice as an answer for the equation of the tangent line? $$(y-2)=(\frac{x(x^2-5)}{2(y^2-2)}) (x-0)$$ Or would i have to solve from there? – soniccool Jun 29 '12 at 23:48
Ah i totally forgot about that Alright awesome! :D – soniccool Jun 29 '12 at 23:56
HINT
Yes, you need to implicitly differentiate. In order to describe a line, you need a point and a slope. You already have the point, so you just need to extract the slope from the implicit differentiation.
-
Awesome okay let me try this – soniccool Jun 29 '12 at 23:04
Differentiate both sides with respect to $x$. On the right, we get $4x^3-10x$. On the left we get $(4y^3-8y)\frac{dy}{dx}$.
Now (crucial) we need to evaluate the derivative at $(0,-2)$. At this point, $4x^3-10x$ is $0$, and $4y^3-8y$ is not. So $\left(\frac{dy}{dx}\right)_{(0,-2)}=0$.
Slope is $0$, we pass through $(0,-2)$, so the tangent line is horizontal, and has equation $y=-2$.
- | 2014-09-15T04:44:29 | {
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https://es.mathworks.com/help/symbolic/numeric-to-symbolic-conversion-1.html | # Numeric to Symbolic Conversion
This topic shows how Symbolic Math Toolbox™ converts numbers into symbolic form. For an overview of symbolic and numeric arithmetic, see Choose Numeric or Symbolic Arithmetic.
To convert numeric input to symbolic form, use the `sym` command. By default, `sym` returns a rational approximation of a numeric expression.
```t = 0.1; sym(t)```
```ans = $\frac{1}{10}$```
`sym` determines that the double-precision value `0.1` approximates the exact symbolic value $\frac{1}{10}$. In general, `sym` tries to correct the round-off error in floating-point inputs to return the exact symbolic form. Specifically, `sym` corrects round-off error in numeric inputs that match the forms $\frac{\mathit{p}}{\mathit{q}}$, $\frac{\mathit{p}\pi }{\mathit{q}}$, ${\left(\frac{\mathit{p}}{\mathit{q}}\right)}^{\frac{1}{2}}$, ${2}^{\mathit{q}}$, and ${10}^{\mathit{q}}$, where $\mathit{p}$ and $\mathit{q}$ are modest-sized integers.
For these forms, demonstrate that `sym` converts floating-point inputs to the exact symbolic form. First, numerically approximate $\frac{1}{7}$, $\pi$, and $\frac{1}{\sqrt{2}}$.
`N = [1/7 pi 1/sqrt(2)]`
```N = 1×3 0.1429 3.1416 0.7071 ```
Convert the numeric approximations to exact symbolic form. `sym` corrects the round-off error.
`S = sym(N)`
```S = $\left(\begin{array}{ccc}\frac{1}{7}& \pi & \frac{\sqrt{2}}{2}\end{array}\right)$```
You can force `sym` to accept the input as is by placing the input in quotes. Demonstrate this behavior on the previous input `0.142857142857143`. The `sym` function does not convert the input to `1/7`.
`sym('0.142857142857143')`
`ans = $0.142857142857143$`
When you convert large numbers, use quotes to exactly represent them. Demonstrate this behavior by comparing `sym(133333333333333333333)` with `sym('133333333333333333333')`.
`sym(1333333333333333333)`
`ans = $1333333333333333248$`
`sym('1333333333333333333')`
`ans = $1333333333333333333$`
You can specify the technique used by `sym` to convert floating-point numbers using the optional second argument, which can be `'f'`, `'r'`, `'e'`, or `'d'`. The default flag is `'r'`, for rational form.
### Conversion to Rational Symbolic Form
Convert input to exact rational form by calling `sym` with the `'r'` flag. This is the default behavior when you call `sym` without flags.
```t = 0.1; sym(t,'r')```
```ans = $\frac{1}{10}$```
### Conversion by Using Floating-Point Expansion
If you call `sym` with the flag `'f'`, `sym` converts double-precision, floating-point numbers to their numeric value by using $\mathit{N}\cdot {2}^{\mathit{e}}$, where $\mathit{N}$ and $\mathit{e}$ are the exponent and mantissa respectively.
Convert `t` by using a floating-point expansion.
`sym(t,'f')`
```ans = $\frac{3602879701896397}{36028797018963968}$```
### Conversion to Rational Symbolic Form with Error Term
If you call `sym` with the flag `'e'`, `sym` returns the rational form of `t` plus the error between the estimated, exact value for `t` and its floating-point representation. This error is expressed in terms of `eps` (the floating-point relative precision).
Convert `t` to symbolic form. Return the error between its estimated symbolic form and its floating-point value.
`sym(t,'e')`
```ans = $\frac{\mathrm{eps}}{40}+\frac{1}{10}$```
The error term `eps/40` is the difference between `sym('0.1')` and `sym(0.1)`.
### Conversion to Decimal Form
If you call `sym` with the flag `'d'`, `sym` returns the decimal expansion of the input. The `digits` function specifies the number of significant digits used. The default value of `digits` is 32.
`sym(t,'d')`
`ans = $0.10000000000000000555111512312578$`
Change the number of significant digits by using `digits`.
```digitsOld = digits(7); sym(t,'d')```
`ans = $0.1$`
For further calculations, restore the old value of `digits`.
`digits(digitsOld)`
### Conversion to Variable Precision
You can create symbolic numbers with variable-precision floating-point arithmetic by using `vpa`. By default, `vpa` calculates values to 32 significant digits.
`piVpa = vpa(pi)`
`piVpa = $3.1415926535897932384626433832795$`
When you use `vpa` on a numeric input, such as `log(2)`, the numeric expression is first evaluated to the MATLAB® default double-precision number that has less than 32 significant digits. Then, `vpa` is applied on that double-precision number, which can be less accurate. For more accurate results, convert numeric expressions to symbolic expressions with `sym` and then use `vpa` to evaluate the results with variable precision. For example, find `log(2)` with 17 and 20 digits precision.
`vpaOnDouble = vpa(log(2))`
`vpaOnDouble = $0.69314718055994528622676398299518$`
`vpaOnSym_17 = vpa(log(sym(2)),17)`
`vpaOnSym_17 = $0.69314718055994531$`
`vpaOnSym_20 = vpa(log(sym(2)),20)`
`vpaOnSym_20 = $0.69314718055994530942$`
When you convert large numbers, use quotes to exactly represent them.
`inaccurateNum = vpa(123456789012345678)`
`inaccurateNum = $123456789012345680.0$`
`accurateNum = vpa('123456789012345678')`
`accurateNum = $123456789012345678.0$` | 2023-03-31T09:07:25 | {
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https://math.stackexchange.com/questions/340631/are-there-countably-or-uncountably-many-infinite-subsets-of-the-positive-even-in | # Are there countably or uncountably many infinite subsets of the positive even integers?
Let $S$ be the set of all infinite subsets of $\mathbb N$ such that $S$ consists only of even numbers.
Is $S$ countable or uncountable?
I know that set $F$ of all finite subsets of $\mathbb N$ is countable but from that I am not able to deduce that $S$ is uncountable since it looks hard to find a bijection between $S$ and $P(\mathbb N)\setminus F$. Also I am not finding the way at the moment to find any bijection between $S$ and $[0,1]$ to show that $S$ is uncountable nor I can find any bijection between $S$ and $\mathbb N$ or $S$ and $\mathbb Q$ to show that it is countable. So I am thinking is there some clever way to show what is the cardinality of $S$ by avoiding bijectivity arguments?
So can you help me?
Notice that by dividing by two, you get all infinite subsets of $\mathbb{N}$. Now to make a bijection from $]0,1]$ to this set, write real numbers in base two, and for each real, get the set of positions of $1$ in de binary expansion.
You have to write numbers of the form $\frac{n}{2^p}$ with infinitely many $1$ digits (they have two binary expansions, one finite, one infinite). Otherwise, the image of such a real by this application would not fall into the set of infinite sequences of integers (it would have only finitely many $1$).
• Note that $0.01111\ldots_2 = 0.1_2$. Of course your argument can be easily fixed, but the function you describe is not a bijection. – dtldarek Mar 25 '13 at 14:02
• That's exactly why I wrote one need to use the infinite expansion when there are two possible binary expansions. Why do you think it's not a bijection? – Jean-Claude Arbaut Mar 25 '13 at 14:04
• You are right, I should refresh my browser more often (sometimes AJAX does not work). – dtldarek Mar 25 '13 at 14:06
• But you are right, I had to include 1 in the interval, to take care of image of $\mathbb{N}$. Now I think everything is fixed. – Jean-Claude Arbaut Mar 25 '13 at 14:08
• Thank you, it surely looks good! – user67878 Mar 25 '13 at 14:08
The set of even numbers is countably infinite, and so equinumerous to $\Bbb N$ (explicit bijection given by $n\mapsto\frac12n$). With $F$ the set of finite subsets of $\Bbb N$, as you say, we can use the map $f(n)=\frac12n$ to get a bijection $g:S\to P(\Bbb N)\setminus F$, given by $$g(A)=\{f(n):n\in A\}.$$ I leave it to you to prove that it's bijective.
• So, if I understood your argument well, it is the case that if there is a bijection between $A$ and $B$ then there is a bijection between $R$ and $S$ where $R$ is the set of all infinite subsets of $A$ and $S$ is the set of all infinite subsets of $B$? – user67878 Mar 25 '13 at 13:57
• Absolutely right, and you can prove that as a general statement using the same idea that I outlined above. – Cameron Buie Mar 25 '13 at 13:59
• Do you mean by the "general statement" that it works for every cardinality of $A$ and $B$ if we precisely enough define cardinality of $R$ and $S$, which is not needed in my case because every infinite subset in my problem must be of countable infinity? – user67878 Mar 25 '13 at 14:05
• What a flaw by me! It is fairly "obvious" what is the set of all infinite subsets of some set so in my previous comment maybe it is better to put away the word "precise". And you already answered that it is possible so it looks like I am asking the same question as in the first comment. It was a wave of confusion few minutes ago so it is clearer now for me to see it should work for every cardinality. Thank you! – user67878 Mar 25 '13 at 14:14
• So, if I ask you what approach would you choose for yourself, a set theory in which that statement is true, or set theory in which that statement is not true, what would you choose? You can choose it without an axiom of choice! ;.D – user67878 Mar 27 '13 at 19:44
Use the Cantor-Bernstein theorem, that is (in the simplified version)
$$|A| \leq |B| \text{ and } |B| \leq |A| \text{ implies } |A| = |B|.$$
We have that $S \subset P(\mathbb{N})$, so $|S| \leq |P(\mathbb{N})|$. Moreover, we have an injection $f : P(\mathbb{N}) \to S$ given by
$$f(A) = \{4k \mid k \in A\} \cup \{4k+2 \mid k \in \mathbb{N}\}.$$
Hence, we have $|P(\mathbb{N})| \leq |S|$, so $|S| = |P(\mathbb{N})|$.
I hope this helps ;-)
• Yes, you all got +1 for the effort! – user67878 Mar 25 '13 at 14:02
• @Thus As there might be many solutions, it would be still of benefit to you to know the Cantor-Bernstein theorem. Sometimes bijections are troublesome to construct (see, e.g. bijection between $P(\mathbb{N})$ and $\mathbb{R}$), while usually there is a much more simple way to construct two injections. – dtldarek Mar 25 '13 at 14:14
• I know about the theorem but didn´t look in that direction that it could be easier to construct two injections instead of one bijection, not because of $1<2$ :-D – user67878 Mar 25 '13 at 14:21
Hint: If $\{x_n\mid n\in\Bbb N\}$ is a set of even integers then $\{\frac{x_n}2\mid n\in\Bbb N\}$ is a sequence of integers. Show that this operation is in fact a bijection between $S$ and the set of all infinite subsets of integers. The latter is more familiar.
• Thank you! I have a question that I can put in a comment and I believe that you could solve it, do you want of me to ask you? – user67878 Mar 25 '13 at 14:28
• If it is related, sure. – Asaf Karagila Mar 25 '13 at 14:29
• I will answer it too, I just want of you to confirm is my line of reasoning correct. Let $S$ be the set of all subsets of $\mathbb N$ such that every set that belongs to $S$ consists only of numbers that are pairwise coprime (relatively prime). So I was thinking of this. The set of all prime numbers $P$ is in $S$ and it is infinite. Because the $P$ is infinite and every infinite subset of $P$ is also obviously consisted of coprime numbers then $S$ is uncountable because the $P$ is countable so it has uncountable number of countable subsets. Is this good reasoning? – user67878 Mar 25 '13 at 14:43
• If it is then I understood your arguments for the question I asked. – user67878 Mar 25 '13 at 14:43
• Yes, it is an excellent reasoning. – Asaf Karagila Mar 25 '13 at 14:48 | 2021-02-26T16:33:38 | {
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http://mooboo.com.br/guardian-meaning-xhxdp/grnhj3.php?tag=scalene-obtuse-triangle-ddef80 | a) The length of one side and the perpendicular distance of that side to the opposite angle. The sum of all the angles in any triangle is 180°. This is because scalene triangles, by definition, lack special properties such as congruent sides or right angles. Calculates the other elements of a scalene triangle from the selected elements. There are two ways to classify triangles: by their sides and their angles, like sails out on the high seas can be right or isosceles. Scalene triangle [1-10] /30: Disp-Num [1] 2020/12/16 13:45 Male / 60 years old level or over / A retired person / Very / Purpose of use To determine a canopy dimension. Reduced equations for equilateral, right and isosceles are below. Obtuse Scalene Triangle Translation to prove SSS Congruence 1. Equilateral triangle. Or look at the foot of this goose; it’s scalene and obtuse. Obtuse Angled Triangle: A triangle whose one of the interior angles is more than 90°. scientific notation A method for writing extremely large or small numbers compactly in which the number is shown as the product of two factors. Scalene Triangle Equations These equations apply to any type of triangle. An obtuse triangle has one angle measuring more than 90º but less than 180º (an obtuse angle). Area of a triangle. The converse of this is also true - If all three angles are different, then the triangle is scalene, and all the sides are different lengths. Questionnaire. carotid triangle, superior carotid trigone. Yes! how to i find the length in a Scalene triangle? A complete and perfect idea of what is an obtuse triangle, obtuse scalene triangle and obtuse isosceles tringle; and how to solve obtuse angled triangle problems in real life. Learn more. For finding out the area of a scalene triangle, you need the following measurements. Less than 90° - all three angles are acute and so the triangle is acute. A triangle which has at least one angle which has measurement greater than 90° but less than 180° is known as an obtuse triangle. Types of Triangle by Angle. [insert scalene G U D with ∠ G = 154° ∠ U = 14.8° ∠ D = 11.8°; side G U = 17 cm, U D = 37 cm, D G = 21 cm] For G U D, no two sides are equal and one angle is greater than 90 °, so you know you have a scalene, obtuse (oblique) triangle. In any triangle, two of the interior angles are always acute (less than 90 degrees) *, so there are three possibilities for the third angle: . A scalene triangle may be right, obtuse, ... scalene triangle A triangle with three unequal sides. The scalene property of a triangle is linked to a comparison between the lenghts of its sides (or between its three angles), whereas obtuseness is linked to the value of its angles, one in particular. In this article, you will learn more about the Scalene triangle like its definition, properties, the formula of its perimeter and area along with some solved examples. Acute triangle. Area of Scalene Triangle Formula. An obtuse-angled triangle is a triangle in which one of the interior angles measures more than 90° degrees. Pythagorean Theorem. b) has 3 sides of equal length. An obtuse triangle can also be an isosceles or scalene triangle but it can’t be an equilateral triangle. In an equiangular triangle, all the angles are equal—each one measures 60 degrees. Most triangles drawn at random would be scalene. It means all the sides of a scalene triangle are unequal and all the three angles are also of different measures. Hypotenuse. Return to the Shape Area section. Perimeter of a triangle. Use details and coordinates to explain how the figure was transformed, including the translation rule you applied to your triangle. The triangles above have one angle greater than 90° Hence, they are called obtuse-angled triangle or simply obtuse triangle.. An obtuse-angled triangle can be scalene or isosceles, but never equilateral. It is not possible to draw a triangle with more than one obtuse angle. Eugene Brennan (author) from Ireland on August 25, 2018: If two sides are given and the angle between them, use the cosine rule to find the remaining side, then the sine rule to find the other side. Perimeter: Semiperimeter: Area: Area: Base: Height: Angle Bisector of side a: Angle Bisector of side b: Find scalene triangle stock images in HD and millions of other royalty-free stock photos, illustrations and vectors in the Shutterstock collection. If c is the length of the longest side, then a 2 + b 2 < c 2, where a and b are the lengths of the other sides. What is Obtuse Triangle? Describe the translation you performed on the original triangle. triangle [tri´ang-g'l] a three-cornered object, figure, or area, such as a delineated area on the surface of the body; called also trigone. Click (but don't drag) the mouse cursor at the first, second, and third corners … Drag the orange vertex to reshape the triangle. Are you bored? The Formula for Scalene Triangle. Obtuse triangle. Pythagorean triples. Was this site helpful? scattergram A graph with points plotted on a coordinate plane. Interior angles are all different. Congruent triangles. A triangle with an interior angle of 180° (and collinear vertices) is degenerate. Here, the triangle ABC is an obtuse triangle, as ∠A measures more than 90 degrees. (ii) Isosceles triangle: If two sides of a triangle are equal, then it is called an isosceles triangle. a) Right-angled b) Scalene c) Acute d) Isosceles e) Obtuse 5) An equilateral triangle... a) has 3 different lengths and 3 different angles. Obtuse triangle definition, a triangle with one obtuse angle. A scalene triangle is a triangle where all sides are unequal. Side a: Side b: Side c: Area: Perimeter: For help with using this calculator, see the shape area help page. Describe the translation you performed on the original triangle. Isosceles triangle. Use details and coordinates to explain how the figure was transformed, including the translation rule you applied to your triangle. It may come in handy. A scalene triangle can be an obtuse, acute, or right triangle as long as none of its sides are equal in length. b) The lengths of all three sides. That triangle would also be called right if a ninety degree angle is inside. carotid triangle, inferior that between the median line of the neck in front, the sternocleidomastoid muscle, and the anterior belly of the omohyoid muscle. Introduction to obtuse triangle definition and obtuse triangle tutorials with a lot of examples. To see why this is so, imagine two angles are the same. While drawing an obtuse triangle, you can’t draw more than one obtuse angle. scalene triangle definition: 1. a triangle with three sides all of different lengths 2. a triangle with three sides all of…. c) has all 3 sides the same length and each inside angle is the same size. A triangle with one interior angle measuring more than 90° is an obtuse triangle or obtuse-angled triangle. with three different sides they’re called scalene. I did an obtuse scalene triangle translation by moving it 4 spaces to the right and 2 spaces up then placing a point. Step 1: Since it is an isosceles triangle it will have two equal angles. In geometry, Scalene Triangle is a triangle that has all its sides of different lengths. Step 2: Let x be one of the two equal angles. Try the Fun Stuff. Scalene Triangle: No sides have equal length No angles are equal. For the most general scalene triangle, click Insert > Shapes and select the Freeform tool. set A well-defined group of objects. Note: It is possible for an obtuse triangle to also be scalene or isosceles. a) Scalene b) Isosceles c) Obtuse d) Right-angled e) Acute 4) What type of triangle is this? The interior angles of a scalene triangle are always all different. Obtuse Scalene Triangle Translation to prove SSS Congruence 1. What is a Scalene Triangle? The angles in the triangle may be an acute, obtuse or right angle. A triangle that has an angle greater than 90° See: Acute Triangle Triangles - Equilateral, Isosceles and Scalene we konw only one angle and one length. Area of Scalene Triangle With Base and Height (i) Equilateral triangle: If all sides of a triangle are equal, then it is called an equilateral triangle. An obtuse triangle is a type of triangle where one of the vertex angles is greater than 90°. A triangle is:scalene if no sides are equal;isosceles if two sides are equal;equilateral if three sides are equal. BookMark Us. Triangle. See more. A scalene triangle may be right, obtuse, or acute (see below). A scalene triangle is one where none of the 3 sides are equal. Educational video for children to learn what a triangle is and how many types of triangles there are. But the triangle you sketch should be a non-right-angle, scalene triangle (as opposed to an isosceles, equilateral, or right triangle). select elements \) Customer Voice. Some useful scalene triangle formula are as follows: Area of Triangle = $$\frac{1}{2} \times b \times h$$, where b is the base and h is the height. Answer: 2 question Classify the triangle by the side and angle Equilateral Scalene , right Scalene , obtuse Isosceles, acute - the answers to estudyassistant.com In an obtuse triangle, if one angle measures more than 90°, then the sum of the remaining two angles is less than 90°. Thousands of new, high-quality pictures added every day. An equiangular triangle is a kind of acute triangle, and is always equilateral. An obtuse triangle is one where one of the angles is greater than 90 degrees. A(6,8) B(-1,4) C(5,4) is the obtuse triangle … FAQ. The given 96º angle cannot be one of the equal pair because a triangle cannot have two obtuse angles. A triangle is a polygon made up of 3 sides and 3 angles.. We can classify triangles according to the length of their sides. This is one of the three types of triangles, based on sides.. We are going to discuss here its definition, formulas for perimeter and area and its properties. Angles measures more than 90° degrees are also of different lengths 2. a triangle with one angle! Are always all different sides all of…, as ∠A measures more than but. If all sides of a triangle with more than 90° degrees pair because a triangle are equal 5,4 is. There are are equal—each one measures 60 degrees, the triangle ABC is obtuse... 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Measuring more than 90° measures 60 degrees triangle ABC is an isosceles or triangle. ) the length in a scalene triangle stock images in HD and of! | 2021-07-23T16:07:23 | {
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http://new.asia-pacific.tv/topcashback/patriot/who/11327792aca6dfcee35719075e78b22 | # partition of sets in discrete mathematics
In other words, the elements of P(X) are subsets of X, and every subset of X is in fact a member of P(X). A representative of the class S. Example: m n (mod d) iff d | (m - n . Sometimes we will call the subsets that make up a partition blocks. If D is the set of cards in a standard 52-card deck, the same-color-as relation on D - which can be denoted ~ C - has two equivalence classes: the sets {red cards} and {black cards}. Note how our definition allows us to partition infinite sets, and to partition a set into an infinite number of subsets. The subsets in a partition are often referred to as blocks. DISCRETE MATHEMATICS ELSEVIER Discrete Mathematics 152 (1996) 47-54 Partitions of graphs into one or two independent sets and cliques Andreas BrandsHidt Universitiit Rostock, FB Informatik, D 18051 Rostock, Germany Received 12 February 1991; revised 14 June 1994 Abstract It is shown in this note that it can be recognized in polynomial time whether the vertex set of a finite undirected graph . The subsets in a partition are often referred to as blocks. Partitions of Sets If X is a set, then the power set of X is the set P(X) consisting of all subsets of X. These types of graphs are known as isomorphism graphs. UMASS AMHERST MATH 300 SP '05, F. HAJIR HOMEWORK 3: SETS AND MAPS 1. Equivalence Relations 3 . Use Venn diagrams to prove set identities 10. Yes, {} is a subset of every set. North East Kingdom's Best Variety best order to read the old testament; sandman hotel victoria bed bugs; yamashiro hollywood parking; charles edward williams obituary; duke dennis discord server link. Discrete Maths: Disjoint Sets | Partitions of SetDiscrete Mathematics playlist: https://www.youtube.com/playlist?list=PL1w8k37X_6L_M7IBbrygh_OPZlpaQ_49a#Part. English. We denote this by aRb. Search: Discrete Math Test 1. A partition of set $$A$$ is a set of one or more nonempty subsets of $$A\text{:}$$ $$A_1, A_2, A_3, \cdots\text{,}$$ such that every element of $$A$$ is in exactly one set. That means two different graphs can have the same number of edges, vertices, and same edges connectivity. (3) Each subset is non-empty. There are two important examples which have their own names: The odd-even topology is the topology where. No, B1 and B3 are not disjoint. Find the union of all A as follows. Let p A ( n , k ) denote the number of multi-color partitions of n into parts in { a 1 , , a k }. The set of all 2x2 matrices with elements from a given set @W is partitioned into a finite number of classes. The set of all 2x2 matrices with elements from a given set @W is partitioned into a finite number of classes. Compute a Frobenius number: Frobenius number {4, 7, 12} Find the partitions of an integer: integer partitions of 10. Williams Syndrome (WS) is a developmental condition that has been shown to have slower development of perceptual . Colour names are used to partition colour space into discrete colour categories. Subjects. MCQ (Multiple Choice Questions with answers about Discrete Mathematics Equivalence Classes Partitions Which of the following is an equivalence relation on R, for a, b ? (1) The union of the subsets is the entire class. Operations can be dened on sets creating an "algebra." Counting the number of elements in a set and counting subsets with a certain property are fundamental in computing probabilities and statistics. Suppose that A is the set of sophomores at your school and B is the set of students in discrete mathematics at your school. Binary relations A (binary) relation R between the sets S and T is a subset of the cartesian product S T. In , a poset P t , t + 1 = 1 i t 1 { x N | ( i 1 ) ( t + 1 ) + 1 x i t 1 } is introduced by Anderson, whose partial order is specified by requiring that x covers y if x y = t or t + 1. Finite and countable sets are fundamental primitives of discrete math- ematics. Symbolically, $$\displaystyle A_1 \cup A_2 \cup A_3 \cup \cdots = A$$ If $$i \neq j$$ then $$A_i \cap A_j = \emptyset$$ There are Q questions that have to be answered. Then we follow the poset method used by Stanley and Zanello to obtain the sum of elements of all ideals I J ( P t , t + 1 ). Symbolically, A1 A2 A3 = A. That means two different graphs can have the same number of edges, vertices, and same edges connectivity. In the quotient ring R = Z / 3 these become equalities: 7 = 4 = 1 = 2 = 5 = 8 (b) If i j then A i A j = . Abstract Euler's identity equates the number of partitions of any non-negative integer n into odd parts and the number of partitions of n into distinct parts. , such that every element of A is in exactly one set. What is partition discrete mathematics? In other words, the elements of P(X) are subsets of X, and every subset of X is in fact a member of P(X). Does "part" mean subset, or does it mean element? Discrete Mathematics Lecture 12 Sets, Functions, and Relations: Part IV 1 . Of course, if $$A$$ is finite the number of subsets can be no larger than $$\lvert A \rvert \text{. We extract the set S 1 of singleton pairs and the set L 1 of left-point pairs (of adjacency pairs) from 0. If B 1 / . 1Set Theory Set Notation and Relations Basic Set Operations Cartesian Products and Power Sets Binary Representation of Positive Integers Summation Notation and Generalizations 2Combinatorics Basic Counting Techniques - The Rule of Products Permutations Partitions of Sets and the Law of Addition Combinations and the Binomial Theorem 3Logic A partition of set A is a set of one or more nonempty subsets of A: A1, A2, A3, , such that every element of A is in exactly one set. Similar observations can be made to the equivalence class {4,8} . 1 The comments suggest that the main terminology you need is congruence modulo n . adventure holidays uk for adults; dreamfall: the longest journey; laal singh chaddha trailer release date; sets in discrete mathematics. Assuming "discrete math" is a general topic | Use as referring to a mathematical definition or a periodical instead. Other files and links. Since there are exactly three parts and elements 1, 2, 3 are in different parts, you may as well call the parts they are each in "Part 1 ", "Part 2 " and "Part 3 " respectively. Definition 2.3.1: Partition. It concerns all integers, i.e, 1 2 5 mod 3 for example. In other words, if the intersection of the sets is the empty set. The 2-part partition corresponding to ~ C has a refinement that yields the same-suit-as relation ~ S, which . The Relation Induced by a Partition. Boolean algebra calculator circuit for calculate the branch of mathematics that the branch of mathematics that involves in manipulating & simplifying the logical Discrete Mathematics, Algebra, Textbook, Curriculum, Electronics, Digital, The calculator works for both numbers and expressions containing variables ####How to use this calculator#### Simply enter integers whose greatest common . Express each of these sets in terms of A and B. partition is finer than the partition given. The set {} is not an element of every set. B1 = { n Z : n = 2k, for some integer K} . [Discrete Math] Partition of 3 sets. A family F 2 X is called partition-free if it has no pairwise disjoint members whose union is X. Denoting the maximum of w p ( F ) over all partition-free families F 2 X by m ( n , p ) we prove the rather surprising fact that while m n , 1 k = 1 1 k for all integers k 2, m ( n , p ) 1 as n for all other values of p. Find the union of the sets as follows. . The subsets in a partition are often referred to as blocks. Partitions are one of the core ideas in discrete mathematics. Since A in example 1 is given by A= {2,4,6,8,10}, we can easily verify. Set theory is the foundation of mathematics. PDF Discrete Mathematics . We can use our partition to define what it means for two students to be equivalent, by saying that two students in the class are equivalent if they have the same class rank. (1,2) . Share. If S = T we say R is a relation on S. For example, one possible partition of {1, 2, 3, 4, 5, 6} is {1, 3}, {2}, {4, 5, 6}. Outline Equivalence Relations Partial Orderings 2 . Explanations. 3.1.1Partitions of Sets Recall that a partition of a set A is a set of subsets of A such that every element of A is in exactly one of the subsets. We examine several arithmeti. 8. Do the sets B1, B2, and B3 form a partition of the universal set and why? Beck conjectured and Andrews proved th. [ P i { } for all 0 < i n ] The union of the subsets must equal the entire original set. A group of 21 students participates in a discrete mathematics competition. / B k is a partition of {1, ., n} with k > 1 blocks, then a connector is an ordered pair . , such that every element of A is in exactly one set. A partition of set A is a set of one or more nonempty subsets of A: A1, A2, A3, , such that every element of A is in exactly one set. In mathematics, a partition of a set is a grouping of its elements into non-empty subsets, in such a way that every element is included in exactly one subset. Consider a relation R on a set S satisfying the following properties: R is antisymmetric, i.e., if xRy and yRx, then x = y. R is transitive, i.e., xRy and yRz, then xRz. = A. Symbolically, A1 A2 A3 = A. Show that the distinct equivalence classes in example 1 form a partition of the set A there. The example of an isomorphism graph is described as follows: . Solution In example 1 we have shown that [2]= {2,6,10} and [4]= {4,8} are the only distinct equivalence classes. Symbolically, (a) A 1 A 2 A 3 . }$$ Example 2.3.2. Basic building block for types of objects in discrete mathematics. The union of the subsets must equal the entire original set. Equivalence Classes of R. The Partition Induced by R on a set A. We call S(k, n) \\newcommand{\\hexbox}[3]{ \\def\\B{\\mathbf{B}} When we write \$$\\lambda=\\lambda_1^{i_1}\\lambda_2^{i_2}\\cdots\\lambda_n^{i_n}\\text{,}\$$ we will assume that . The example of an isomorphism graph is described as follows: sets in discrete mathematics December 9, 2021. Abstract Euler's identity equates the number of partitions of any non-negative integer n into odd parts and the number of partitions of n into distinct parts. then R is an equivalence relation, and the distinct equivalence classes of R form the original partition {A 1, ,A n}.. Found inside - Page 92Find the number of subsets of X such that each subset has r elements and no two elements in a subset are consecutive integers. Definition 2.3.1: Partition. Discrete Mathematics and Its Applications Kenneth Rosen. We call the subsets that make up the partition blocks or parts of the partition. Discrete Mathematics and Combinatorics; Access to Document. Set operations in programming languages: Issues about data structures used to represent sets and the computational cost of set operations. If i j then Ai Aj = . Abstract. Many different systems of axioms have been proposed. In particular, we find explicit formulas for the total perimeter and the total site-perimeter over all set partitions of [n]. Examples for Discrete Mathematics. 10.1016/j.ejc.2009.07.001. Abstract. generate all partitions of a set (7) . In mathematics, the partition topology is a topology that can be induced on any set X by partitioning X into disjoint subsets P; these subsets form the basis for the topology. Discrete Mathematics 1. Construct partition such that sum of chromatic numbers is greater than chromatic number of graph Sh. If (a,b) R, we say a is in relation R to be b. The principal object of this paper is to estimate how small a matrix is guaranteed to contain an rxs submatrix all of whose 2x2 submatrices . In this paper, we study the generating function for the number of set partitions of [n] represented as bargraphs according to the perimeter/site-perimeter. CS 441 Discrete mathematics for CS M. Hauskrecht Set difference Definition: Let A and B be sets. Determine the power set of a set General denitions: set A collection of discrete items, whether numbers, letters, people, animals, cars, atoms, planets, etc. A set of n elements can be partitioned into k unordered subsets of r elements each (kr = n) in the following number of ways: 1 k! Learn vocabulary, terms, and more with flashcards, games, and other study tools. Discrete math sets, relations, functions, and graphsour experts know it all In mathematics you will often encounter statements of the form "A if and only if B" or "A $\Leftrightarrow$ B" Discrete Mathematics Online Lecture Notes via Web With its clear presentation, the text shows students how to present cases logically beyond this course . Hence the number 3 has 3 partitions: 3 2+1 1+1+1 The number of partitions of n is given by the partition function p ( n ). You'll learn how to count the number of ways to partiti. Enumeration of Gap-Bounded Set Partitions. German mathematician G. Cantor introduced the concept of sets. Determine whether sets form a partition of a given set 11. How do we count set partitions? The isomorphism graph can be described as a graph in which a single graph can have more than one form. Link to publication in Scopus. Let A be a set with a partition and let R be the relation induced by a partition, then R is reflexive, symmetric, and transitive. For a set of the form A = {1, 2, 3, ., n}.It is called partition of the set A, a set of k<=n elements which respect the following theorems:. This book will help you think well about discrete problems: problems like chess, in which the moves you make are exact, problems where tools like calculus fail because there's no continuity, problems that appear all the time in games, puzzles . The median m that partitions the scores into two equal-sized groups: Those below m and those above m. Zermelo-Fraenkel set theory (ZF) is standard. In a similar way, we can derive a formula for the number of unordered partitions of a set. The set {} is a subset of every set. So we have a b mod 3 3 a b in the ring Z. The subsets in a partition are often referred to as blocks. the money for below as .
3. P n = S ] The intersection of any two distinct sets is empty. Each of the remaining 100 3 = 97 parts can be in any of these three parts, meaning that there are 3 97 partitions which meet your conditions. The set S is called the domain of the relation and the set T the codomain. . Combinatorics Compute binomial coefficients (combinations): 30 choose 18. Recall that a partition of a set S is a collection of mutually disjoint subsets of S whose union is all of S. In other words, every element of S belongs to exactly one of the subsets of the partition. Let 1 be the remaining partition of the set [ n] ( S 1 L 1). This video explains set partitions and the combinatorics behind them. Symbolically, (a) A 1 A 2 A 3 . Beck conjectured and Andrews proved th. Here A 1 = { 1, 2 }, A 2 = { 3, 4 }, A 3 = { 5, 6 } .
#### partition of sets in discrete mathematics
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※アジア太平洋大家の会無料メルマガをお送りします。 | 2022-09-25T21:17:54 | {
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https://www.physicsforums.com/threads/discontinous-functions-question.645832/ | # Homework Help: Discontinous functions question.
1. Oct 21, 2012
### peripatein
Hi,
For a function to be continuous, three conditions must be met - the function must be defined at a point x0, its limit must exist at that point, and the limit of the function as x approaches x0 must be equal to the value of the function at x0.
Now, assuming my function is f(x)=[x], which assigns an integer smaller than x to f(x). Thus, if 1≤x<2, f(x)=1.
The instructor claimed that this function is discontinuous for every integer x, which is perfectly clear just from looking at the graph of f(x), which is a step graph. He also mentioned that two of the three requirements for continuity are unmet in this case.
Which two requirements of the three above are unmet in this case?
2. Oct 21, 2012
### Staff: Mentor
For the limit to exist, the left- and right-side limits have to exist and be equal. You're dealing with a step function that has jumps at integer values. Take a look at the three conditions at these integer values.
3. Oct 21, 2012
### peripatein
The instructor claimed that for a limit to exist at a point x0, the limit has to be equal both from above and below, unless the limit from one direction does not exist and then the limit at x0 DOES exist. Hence, that requirement is met for the step function as described above, I believe. Is it not?
Another requirement that is met is that the function is defined for any x0. Is it not?
So two requirements are already met and that alone disagrees with the instructor's claim that two conditions are unmet.
Could someone please clarify? What am I missing?
4. Oct 21, 2012
### Staff: Mentor
$$\lim_{x \to 1^-} f(x)?$$
and
$$\lim_{x \to 1^+} f(x)?$$
Do both exist? If so, are they equal?
5. Oct 21, 2012
### Ray Vickson
No, the function f(x) = [x] does not have a limit as x → n, for every integer n. The limits exist in BOTH directions, however. (Why do these two statements not contradict each other?)
RGV
6. Oct 21, 2012
### peripatein
Mark, is the limit from below equal to 0, whereas the limit from above is equal to 1?
7. Oct 21, 2012
### peripatein
Do both limits indeed exist and whereas the first is equal to 0, the second is equal to 1?
8. Oct 21, 2012
### Ray Vickson
What is f(0.9)? What is f(0.99)? What is f(1.1)? What is f(1.01)? Can you see now what is happening?
RGV
9. Oct 21, 2012
### peripatein
Well, as x approaches 1 from below f(x) is still 0 (or does that limit not exist?). When x approaches 1 from above, f(x) is 1. Is that correct?
10. Oct 21, 2012
### Staff: Mentor
Yes. So
$$\lim_{x \to 1^-} f(x) = 0$$
and
$$\lim_{x \to 1^+} f(x) = 1$$
So does $\lim_{x \to 1} f(x)$ exist? | 2018-09-22T05:17:05 | {
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https://math.stackexchange.com/questions/1377728/two-circles-inside-a-semi-circle/1377747 | # Two circles inside a semi-circle
Two circles of radius 8 are placed inside a semi-circle of radius 25.The two circles are each tangent to the diameter and to the semi-circle.If the distance between the centers of the two circles is $\lambda$,then the value of $\frac{60}{\lambda}$.
My try:I am not sure if i am correct,i guess these two circles touch each other and $\lambda$ is 16.Can someone explain me correct answer and correct method.
• Please also share the effort you have put in. – Shailesh Jul 29 '15 at 8:07
• The question did not say they are touching each other. Therefore, $\lamda$ is not 16. – Mick Jul 29 '15 at 8:09
• @Mick,yes question does not say they touch each other. – Vinod Kumar Punia Jul 29 '15 at 8:10
• I am sorry, what is the question? – nicks Jul 29 '15 at 9:20
Try to figure that out from the attached diagram.
• ,you drew this diagram in GeoGebra? – Vinod Kumar Punia Jul 29 '15 at 8:27
• Definitely a +1 for the diagram. A pic is worth a thousand words. I had a hard time describing my answer without being able to draw. – Shailesh Jul 29 '15 at 8:33
• @VinodKumarPunia Yes. It is a great tool. – Mick Jul 29 '15 at 11:37
• @Shailesh In the field of Geogebra, I am a newbie also because many useful applications are still "hidden" (to me). – Mick Jul 29 '15 at 11:41
• @Mick. Good job. Only point is to actually show that the line through centers also passes through the point of tangency. – Shailesh Jul 29 '15 at 13:42
Considering a line from the centre of the semicircle through the centre of one of the circles, the distance from the centre of the semicircle to the centre of the small circle is $25-8=17$. Then applying pythagoras gives the distance between the centres $=\lambda=2\times15=30$
Hint :
From the midpoint of the larger semicircle, draw a line through the center of one of the smaller semi-circle, extending upto the the larger semi-circle. This is the point of tangency. Why ? . Now, let this line make an angle $\theta$, set up equations, and you will find a nice 8-15-17 right triangle. So the 'horizontal' length between the centers is 15. Twice that is 30. So $\lambda$ = 2. I can see that David Quinn has given a similar answer.
Consider the segment with endpoints at the center of the semicircle and at the center of one of the circles. It's length is $25-8=17$. Now, think of a right triangle whit this segment as hypotenuse and with one side upon the diameter. This side has length 15, because of the pythagorean theorem and the fact that the other sides are 8 and 17. The same is true for the other circle and if we add the horizontal distances (consider the diameter of the semicircle is horizontal) we get that $\lambda = 30$
Edit: later i saw that other people posted similar answers | 2020-08-13T21:57:07 | {
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http://math.stackexchange.com/questions/57471/a-n-c-alphan-where-c-in-r-a-general-solution-to-difference-equation-a/57499 | $a_{n}=C \alpha^{n}$ where $C\in R$ a general solution to difference equation $a_{n+1}-\alpha a_{n}=0$?
Suppose $T(N) = 2 T(\frac{N}{2}) + N$, where $T(1) = 0$ and $N = 2^{n}$ where $n \in N$.
I get $\alpha ^{n-1} (\alpha -2) = \frac{n}{C}$, using the general difference equation formula.
My sketch:
3^{n-1} T(n) n
1 0 1
3 2 2
9 8 3
27 24 4
81 62 5
So my guess $O(T(N)) = 3^{N}$.
Now $3^{n-1} (3 -2) = \frac{n}{C}$ so $3^{n-1} = \frac{n}{C}$. So $(n-1) \ln(3) = \ln(n) - \ln(C)$ so
$n - \ln(n) = K$,
where $K=\ln(3) - \ln(C)$.
But I cannot see why it would be a solution.
-
Please make the body of your question self-contained. The title is not part of the message, it's an indexing feature. Books don't ask you to read the spine in order to understand the first page. – Arturo Magidin Aug 14 '11 at 22:34
Either I am missing something or the question and its title are 100% unrelated. – Did Aug 14 '11 at 23:05
Showing your calculations is good, but you never quite said what it is that you would like to show. In your table, the label should not be $T(n)$, you are listing $T(\log_2 N)$. By the way, there is a little error, $62$ should be $64$. If you make the question clear, I can give a sample solution. – André Nicolas Aug 15 '11 at 1:17
We approach the problem in various related ways. First we set up some notation. We are working with powers of $2$. Thus it may be a good idea to let $N=2^n$. So $n=\log_2(N)$. Let $$W(n)=T(N)=T(2^n).$$
Then our recurrence can be rewritten as $$W(n)=2W(n-1)+2^n$$ with the initial condition $W(n)=0$.
First approach: Make a table very much like the table that you drew up. The indices in your table are not quite right, since $T(1)=0$, but $N=1$ corresponds to $n=0$.
Compute. We have $W(0)=0$, $W(1)=2$, $W(2)=8$, $W(3)=24$, $W(4)=64$, $W(5)=160$, $W(6)=384$.
Note that $W(0)=0\cdot 2^0$, $W(1)=2\cdot 2^1$, $W(2)=2\cdot 2^2$, $W(3)=3\cdot 2^3$, and $W(5)=5 \cdot 2^5$. It is reasonable to conjecture that $W(n)=n2^n$.
This is easy to verify by induction. The result holds at $n=0$. Suppose that it holds at $n=k$. We show that it holds at $n=k+1$.
This is straightforward:
$$W(k+1)=2W(k)+2^{k+1}=(2k)(2^k)+2^{k+1}k2^{k+1}+2^{k+1}=(k+1)2^{k+1}.$$
Now let's put things back in terms of $T$ and $N$. Recall that $n=\log_2 N$. Then $$T(N)=W{2^n}=n2^n =N\log_2 N.$$ (At the end, we reversed order of typographical reasons.)
Thus we can write $T(N)=O(N\log_2 N)$. (We have proved something quite a bit stronger.)
Second approach: This is a variant of the first approach. You drew up a table. That was a very good idea. Imagine making a careful table that gives $N$, $n$, and $T(N)$. The pattern is likely to jump out.
Even better, make a simple table that gives $N$, maybe $n$, and definitely $T(N)$, and (important) $T(N)/N$. Immediately, everything becomes clear. Do it!
Unfortunately, the $n$ in your table were off by $1$, $N$ was not given, and there was a little mistake, $62$ should have been $64$. That made it difficult to see a pattern.
Third approach: In the first approach, we arrived at the linear recurrence $W(n)=2W(n-1)+2^n$. Unfortunately the recurrence is not homogeneous. Perhaps you have been taught general techniques for dealing with such recurrences. The homogeneous part $W(n)=2W(n-1)$ is easy to solve. But we also need a particular solution of the full non-homogeneous equation. You may have been taught some techniques, or that may have been dealt with by guessing, which really leaves us no better off than with the first approach.
However, there is a little trick that will make everything collapse! We make this recurrence look nicer by letting $W(n)=2^nQ(n)$. Then our recurrence can be rewritten as $$2^nQ(n)=(2)(2^{n-1}Q(n-1)+2^n,$$ which simplifies to $Q(n)=Q(n-1)+1$, with $Q(0)=0$. The solution of this recurrence is hard to miss!
Fourth approach: We were "lucky" that the recurrence had a simple explicit solution. But even if we had something a little messier, we could guess a simple function $f(N)$ such that $T(N)=O(f(N))$. Now we need to prove it. We want to show there is a constant $k$ such that $T(N) \le kf(N)$.
Then we can use the recurrence to prove that if $T(N)\le kf(N)$, then $T(2N)\le kf(2N)$. In our case, pick $f(N)=N\log N)$. The rest is not difficult. This sort of approach is important, since in general we cannot reasonably expect exact solutions. But until one acquires some experience, it may feel harder than exact solutions (in fact it is easier).
-
I don't understand the fourth approach. Can you really find the formula with somehow testing just upper bounds? How does it work? – hhh Aug 15 '11 at 5:12
@hhh: There is no universal procedure! But let's guess that $T(N) \le 10N\log N$. Now we show that if $T(M) \le 10M\log M)$, then $T(2M)\le 10(2M)\log(2M)$, which settles the matter. The idea is that by the recurrence, $T(2M)=2T(M)+2M$. Use the inequality for $T(M)$. We get $T(2M) \le 20 M\log M +2M$. But $\log(2M)=\log 2 +\log M$, so $T(2M)\le 10(2M)\log(2M) +(2M-20M\log 2)$. Finally, $2M-20M\log 2$ is negative, and we are finished. Not much fun, but routine after the first few times. You have worked a lot more with equations than inequalities, that's why it may seem hard. – André Nicolas Aug 15 '11 at 5:56
You appear to have defined $T$ by the recurrence $T(2^n) = 2T(2^{n-1})+2^n$ for $n>0$, with initial condition $T(1)=0$. At least this more or less matches your table, except that $62$ should be $64$ (as André noted), and your $T(n)$ column is actually $T(2^{n-1})$. To see more clearly what’s going on, replace $T$ by the function $t$ defined as follows: $t(0)=0$, and $t(n)=2t(n-1)+2^n$ for $n>0$. It’s easy to check that $t(n) = T(2^n)$ for all $n \ge 0$. A closed form for $t(n)$ can be found by a variety of elementary means, but perhaps the most efficient is to let $s(n) = \frac{t(n)}{2^n}$; then $s$ satisfies the recurrence $$s(n) = \frac{t(n)}{2^n} = \frac{2t(n-1)+2^n}{2^n} = s(n-1)+1$$ for $n>0$ with initial condition $s(0) = 0$. Clearly $s(n) = n$ for all $n\ge 0$, so $T(2^n) = t(n) = n2^n$ for all $n\ge 0$.
But I cannot see what this has to do with your title or with the equation $\alpha^{n-1}(\alpha-2)=\frac{n}{C}$.
- | 2014-08-01T11:56:30 | {
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https://math.stackexchange.com/questions/1098587/maximum-number-of-edges-in-a-bipartite-graph | # Maximum number of edges in a bipartite graph
Prove that for a bipartite graph $G$ on $n$ vertices the number of edges in $G$ is at most $\frac{n^2}{4}$.
I used induction on $n$.
Induction hypothesis: Suppose for a bipartite graph with less than $n$ vertices the result holds true.
Now take a bipartite graph on $n$ vertices.Let $x,y$ be two vertices in $G$ where an edge exist between $x$ and $y$. Now remove these two vertices from $G$ and consider this graph $G'$. $G'$ has at most ${(n-2)^2}\over4$. Add these two vertices back. Then the number of edges $G$ can have is at most
$$|E(G')|+d(x)+d(y)-1$$
My question is in my proof I took $d(x) + d(y) \le n$, where $d(x)$ denotes the degree of vertex $x$. Can I consider $d(x)+d(y) \leq n$? I thought the maximum number of edges is obtained at the situation $K_{\frac n 2,\frac n 2}$
• Yes- $\frac{n}{2} \cdot \frac{n}{2} = \frac{n^{2}}{4}$. If $n$ is even, this works out very well. Otherwise, $n$ is odd. So one partition will have an extra vertex. – ml0105 Jan 10 '15 at 6:26
• $d(x)+d(y)\leq n$ looks valid. You didn't mention your two base cases to get the induction started. Can you add that? Also a brief nod in the direction of the emptying of one part of G in the case where x or y is the last node in its part. – Joffan Jan 10 '15 at 6:35
There is no need to use induction here. A bipartite graph is divided into two pieces, say of size $p$ and $q$, where $p+q=n$. Then the maximum number of edges is $pq$. Using calculus we can deduce that this product is maximal when $p=q$, in which case it is equal to $n^2/4$.
To show the product is maximal when $p=q$, set $q=n-p$. Then we are trying to maximize $f(p)=p(n-p)$ on $[0,n]$. We have $f'(p)=n-2p$, and this is zero when $p=n/2$. After checking the end points we conclude that the maximum is $n^2/4$ at $p=n/2$.
• How can it be shown that product is maximal when p=q – sam_rox Jan 10 '15 at 7:06
• @sam_rox please see edit. – Matt Samuel Jan 10 '15 at 7:10
• you can also write $p$ as $\frac{n}{2}-l$ and use the trick at the bottom of my result. – Jorge Fernández Hidalgo Jan 10 '15 at 7:20
The sum of the degrees of vertices $x$ and $y$ is indeed less than or equal to $n$. One reason for this is that if a vertex $v$ is adjacent to $x$ it cannot be adjacent to $y$ since $y$ and $v$ would be in the same part.
I would like to share my proof, please tell me what you think of it: The bipartite graph on $n$ vertices with the most edges is clearly a complete bipartite graph, otherwise we could take that graph, add an edge and we would have a graph with more edges.
In a complete bipartite graph with parts of size $k,n-k$ the vertices in the part with $k$ vertices have order $n-k$ and the vertices in the part with $n-k$ vertices have order $k$. The sum of the degrees is then $2k(n-k)$, so by the handshaking lemma the number of edges is $k(n-k)$
Now write $k$ as $\frac{n}{2}-j$. Then $k(n-k)=(\frac{n}{2}-j)(\frac{n}{2}+j)=\frac{n^2}{4}-l^2\leq \frac{n^2}{4}$
suppose the bipartite graph has m vertices in set1 and n vertices in set2 . the sum of degrees of all vertices would be 2 * e= 2 * m * n thus, e = m * n and m+n = v so n*(v-n)= e solve this quadratic equation and discriminant .
Suppose $p,q$ are nonnegative integers with $p+q=n,$ and that $K_{p,q}$ has the maximum number of edges among all bipartite graphs with $n$ vertices.
If we move one vertex from the side with $p$ vertices to the side with $q$ vertices, we lose $q$ edges and gain $p-1$ new edges. Since the number of edges was already maximized, we must have $p-1\le q,$ that is, $p-q\le1.$ Symmetrically, we also have $q-p\le1,$ and so $|p-q|\le1.$
If $n$ is even, then from $p+q=n$ and $|p-q|\le1$ it follows that $p=q=\frac n2,$ and $pq=\frac{n^2}4=\left\lfloor\frac{n^2}4\right\rfloor.$
If $n$ is odd, then $\{p,q\}=\left\{\frac{n-1}2,\frac{n+1}2\right\},$ and $pq=\frac{n^2-1}4=\left\lfloor\frac{n^2}4\right\rfloor.$ | 2019-12-14T13:55:05 | {
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https://mathematica.stackexchange.com/questions/237401/take-function-as-input-for-other-function-access-variable | # Take function as input for other function, access variable
I am very new to Mathematica, and I searched the documentation and google now for quite some time.
Let's say, I want to implement a mathematica function, that does the following mapping: $$f(t)\mapsto \int_0^t f(s)\ \mathrm{d}s+f'(t)$$. ($$f$$ is assumed to be differentiable and integrable.)
Edit: It should also work for $$\mathbb{R}\longrightarrow\mathbb{R}^n$$ functions.
How would you do that? It should be a function like
IntAndDiff[f_]=Integrate[f,{s,0,t}]+D[f,t]
which does of course not work. How am I able to "access" the variable of $$f$$ in this case, so that I can define the integration and differentiation accordingly?
Edit: I also want this function to behave like an operator, so the output should be a $$\mathbb{R}\longrightarrow\mathbb{R}^n$$ function again. I.e. I want to be able to apply this (and other such operators) on the result again...
• ClearAll[IntAndDiff];IntAndDiff[f_] = Integrate[f@s, {s, 0, t}] + D[f@t, t]? – kglr Dec 31 '20 at 0:07
• As @kglr suggested, except use SetDelayed, i.e., IntAndDiff[f_] := Integrate[f[s], {s, 0, t}] + D[f[t], t] Examples: IntAndDiff /@ {Sin, Sin[#] &, Cos, Cos[#] &, Sqrt, Sqrt[#] &, #^2 + 2 # - 3 &} – Bob Hanlon Dec 31 '20 at 0:35
$$f(t)\mapsto \int_0^t f(s)\ \mathrm{d}s+f'(t)$$
A little analysis:
Input is actually a function $$f(t)$$, the output is also a function $$F(t)=\int_0^t f(s)\ \mathrm{d}s+f'(t)$$.
So an easy way is to introduce $$t$$
F[f_, t_]:= Integrate[f[s], {s, 0, t}] + f'[t]
F[Sin,t] (*1*)
Then use Function to make it a pure function
IntAddDiff[f_] := Function[t, Evaluate[Integrate[f[s], {s, 0, t}] + f'[t]]]
Then use # and & to remove t
IntAddDiff[f_] := Evaluate[Integrate[f[s], {s, 0, #}] + f'[#]] &
IntAddDiff[Sin] (*1 &*)
IntAddDiff[#^2 &] (*2 #1 + #1^3/3 &*)
• For me this looks like the solution that works best for my application. I will do some testing, and after this accept and mark the question as answered. – NG98 Jan 3 at 10:29
By using Through we can deal with multiple function {F,G,H} map at t
intAndDiff[f___][t_] :=
Integrate[Through[{f}@s], {s, 0, t}] +
Through[(Derivative[1] /@ {f})@t];
intAndDiff[Sin, Cos, #^2 &][x]
{1, 0, 2 x + x^3/3}
Original
IntAndDiff[f_][t_] := Integrate[f[s], {s, 0, t}] + D[f[t], t];
IntAndDiff[#^2 &][x]
IntAndDiff[Sin][x]
2 x + x^3/3
1
• I think this suggestion is very close to the solution I search for. It does however not work as intended in the case, where f is a function, that takes a scalar as input and returns a vector, as I cannot "concatenate" it. I.e. I search for an operator, so I take a R->R^n function as input and get a R->R^n function as output... – NG98 Jan 2 at 10:34
Same as @cvgmt's solution but defined as an operator, so that it also works on $$\mathbb{R}\to\mathbb{R}^n$$ functions. Both the input and the output are now pure functions:
IntAndDiff[f_] := Function[t, Evaluate[Integrate[f[s], {s, 0, t}] + D[f[t], t]]]
IntAndDiff[#^2 &]
(* Function[t$$, 2 t$$ + t$^3/3] *) IntAndDiff[Sin] (* Function[t$, 1] *)
IntAndDiff[{#, #^2, #^3} &]
(* Function[t$$, {1 + t^2/2, 2 t + t^3/3, 3 t^2 + t$$^4/4}] *) | 2021-04-23T06:10:41 | {
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https://math.stackexchange.com/questions/739804/how-to-understand-cauchys-proof-of-am-gm-inequalitythe-last-step | # How to understand Cauchy's proof of AM-GM inequality(the last step)
The AM-GM inequality: $$a_1a_2\cdots a_n\leq\left(\frac{a_1+\cdots + a_n}{n}\right)^n$$
the trivial case: $a_1a_2 \leq \left(\frac{a_1+a_2}{2}\right)^2$ is self-evident.
then cauchy use this fact repeatedly. He get: $$a_1a_2\cdots a_{2^m} \leq \left(\frac{a_1+a_2+\cdots +a_{2^m}}{2^m}\right)^{2^m}.\tag{\ast}$$
This is easy to understand, and natural, but next step:
he let $$b_1 = a_1,\ b_2=a_2,\ \ldots,\ b_n=a_n,\ b_{n+1}=\cdots=b_{2^m}=\frac{a_1+a_2+\cdots a_n}{n}$$
replace $a_1, a_2,\ldots ,a_{2^m}$ with $b_1, b_2, \ldots b_{2^m}$ in $(\ast)$ and simplified both left and right. He get the answer. I can repeat his proof, but I don't know why Cauchy can got this idea. How does he got it?
• I used to know this... – leo Apr 4 '14 at 17:31
You proved that $$a_1a_2\cdots a_{2^m} \leq \left(\frac{a_1+a_2+\cdots +a_{2^m}}{2^m}\right)^{2^m}\qquad\forall m\in \Bbb N, \tag{\star\star}$$ and you want to prove $$a_1a_2\cdots a_n\leq\left(\frac{a_1+\cdots + a_n}{n}\right)^n\qquad\forall n\in\Bbb N\tag{\star}.$$
You have stablished $(\star)$ for $n=2$ so you have to prove it now for all $n\gt 2$. So pick $n\gt 2$ and $n$ numbers $a_1,\ldots,a_n$.
Now, you can use $(\star\star)$ for a list of numbers of size of power of two. You don't know if $n$ is a power of 2 but you know that sure there is a $m$ so that $2^m\gt n$, so let's add some numbers to our list so we end with a list of size $2^m$.
Which numbers should we add? We don't know, but the most simple thing we can do is add the same number several times (after all this is a try, we don't know if it will work), then our list will be $a_1,\ldots,a_n,\underbrace{x,\ldots,x}_{2^m-n}$.
Using $(\star\star)$ we get $$a_1\cdots a_n\cdot x^{2^m-n}\leq \left(\frac{a_1+\cdots+(2^m-n)x}{2^m}\right)^{2^m}.\tag{\ast}$$
Now you wonder, it would be nice if $x$ had the ability of turn$\frac{a_1+\cdots+a_n+(2^m-n)x}{2^m}$ into $\frac{a_1+\cdots+a_n}{n}$,is there such $x$? That question is answered by the equation $$\frac{a_1+\cdots+a_n+(2^m-n)x}{2^m}=\frac{a_1+\cdots+a_n}{n},$$ which by solving for $x$ tell you that $$x=\frac{a_1+\cdots+a_n}{n}$$ work.
And it turns out that this $x$ does the magic to turn $(\ast)$ into $(\star\star)$. I don't know if Cauchy got it this way, but it's a possible path.
I don't know what Cauchy was thinking, but let me throw some ideas at each other.
When you're trying to prove AM/GM, one natural idea is to think of operations on sequences that preserve their arithmetic mean and to see what they do to the geometric mean (or vice versa). (See Dijkstra's proof, for example.) If you're also thinking you might prove it by induction, you might especially consider operations that add one term to the sequence, or take one away. Putting those two ideas together, you notice that the arithmetic means of $a_1,\dotsc,a_n$ and of $a_1,\dotsc,a_n,\frac{a_1+\dots+a_n}{n}$ are obviously equal. Then you write down the AM/GM inequality for the extended sequence, do some simplifications, and you find that you've proved $$\text{AM/GM for a_1,\dotsc,a_n,\tfrac{a_1+\dots+a_n}{n}} \iff \text{AM/GM for a_1,\dotsc,a_n}$$ and of course, by induction, $$\text{AM/GM for a_1,\dotsc,a_n,\underbrace{\tfrac{a_1+\dots+a_n}{n},\dotsc,\tfrac{a_1+\dots+a_n}{n}}_{m \text{ terms}}} \iff \text{AM/GM for a_1,\dotsc,a_n}$$ At this point you become slightly annoyed because this is exactly what you'd need if your induction went backwards. But then you realize...
His idea is to prove it by mathematical induction, albeit in a different form. Let $\{P(n): n \in \mathbb{N}\}$ be a sequence of propositions.
Key idea: If $P(n)$ is true for infinitely many $n$ and $P(n) \implies P(n-1)$ for all $n$, then $P(n)$ is true for all $n$.
Let $P(n)$ be the A.M-G.M. inequality for $n$ numbers. Cauchy proves it for infinitely many $n$, namely the powers of 2. He then proves the implication $P(n) \implies P(n-1)$.
I don't know how Cauchy found it - perhaps he first noticed that it is easy to prove $P(n) \implies P(n-1)$ than to prove $P(n-1) \implies P(n)$.
• That proof is different from the one OP is talking about. – leo Apr 4 '14 at 17:52
• The proof you mention use this kind of induction. – leo Apr 4 '14 at 19:02 | 2019-05-26T09:47:25 | {
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https://math.stackexchange.com/questions/713112/complex-number-multiplication | Complex number multiplication
$z_1= \cos(4\pi/3) + i\sin(4\pi/3)$
$z_2= \cos(π/3) + i\sin (π/3)$
I want to find out $z_1z_2$.
I know that $(x +iy)(u + iv$) = $(xu - yv) + i(xu + yv)$
So I want to simplify $\cos(4π/3)\cos(π/3) - \sin(4π/3)\sin(π/3)] + i[\cos(4π/3)\cos(π/3) + \sin(4π/3)\sin(π/3)]$
First of, is my formula correct and second, how would I multiply sin or cos functions?
Any assistance would be greatly appreciated.
• Do you know about the polar form of complex numbers? – DonAntonio Mar 15 '14 at 11:57
Well you can simplify using the trigonometric identities $\cos(A \pm B) = \cos A\cos B \mp \sin A \sin B \;\;$ and $\;\; \sin A \cos B \pm \cos A \sin B = \sin (A \pm B)$. For more trigonometric identities check this out.
So the expression reduces to $\cos \left({\frac {5\pi}{3}}\right) + i\sin \left( {\frac{ 5\pi}{3}}\right)$.
But these relationships have been generalised to apply to all multiplications of complex numbers represented in this form. Check this out.
• Your answer is wrong. OP made a multiplication error. You didn't check and continued on that error. – Guy Mar 15 '14 at 12:05
• Correct answer is $\cos (\frac {5\pi}{3}) + i\sin (\frac{5\pi}{3})$ – Guy Mar 15 '14 at 12:05
• @Sabyasachi: Totally my bad. Won't happen again. Edited. – Ishfaaq Mar 15 '14 at 15:12
• And my downvote becomes an upvote :) – Guy Mar 15 '14 at 15:15
Hints:
$$\cos(a\pm b)=\cos a\cos b\mp\sin a\sin b$$
But if you know the polar form, then it is way easier:
$$z_1=\cos\frac{4\pi}3+i\sin\frac{4\pi}3=e^{i\frac{4\pi}3}\;,\;\;z_2=e^{i\frac\pi3}\implies$$
$$z_1z_2=e^{\frac{5\pi}3i}$$
Or even in rectangular form:
$$z_1=-\frac12\left(1+\sqrt3\,i\right)\;,\;\;z_2=\frac12\left(1+\sqrt3\,i\right)\;\ldots$$
• You should probably point out he made a sign error in the multiplication. – Guy Mar 15 '14 at 12:02
• I don't usually read the development of questions when not written with LaTeX...and this wasn't when I first read it. – DonAntonio Mar 15 '14 at 12:03
• Yeah I know, I edited it then. – Guy Mar 15 '14 at 12:04 | 2019-11-15T07:17:23 | {
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https://math.stackexchange.com/questions/3115883/probability-that-displaystyle-vert-x-vert-vert-y-vert-vert-z-vert-vert | # Probability that $\displaystyle \vert x\vert +\vert y\vert +\vert z\vert +\vert x+y+z\vert=\vert x+y\vert +\vert x+z\vert +\vert y+z\vert$
Real numbers $$x, y$$, and $$z$$ are chosen from the interval $$[−1, 1]$$ independently and uniformly at random. What is the probability that $$\vert x\vert +\vert y\vert +\vert z\vert +\vert x+y+z\vert=\vert x+y\vert +\vert x+z\vert +\vert y+z\vert$$
Now if all of $$x, y, z$$ are positive or all negative then the equation is of course satisfied. Hence if we consider a 3D space , it denotes two unit cubes, one in the first octant centred at $$\left( \frac 12,\frac 12,\frac 12\right)$$ and the other in seventh octant centred at $$\left( -\frac 12,-\frac 12,-\frac 12\right)$$.
The total measure of universal set is the cube with edge length $$2$$ centred at origin.
But now I have a problem about what if any two of $$x, y, z$$ are positive while the other remaining be negative or the other way around. Even if I try to make cases it seems to be quite a cumbersome task to approach since we will also need to check signs of $$\vert x+y\vert$$ and similarly others as well as that of $$\vert x+y+z\vert$$
I also thought to give a shot using vectors but didn't reach any specific result.
Any help would be quite beneficial.
Edit:
I would also be happy to see a geometrical intuitive way to attack the problem.
• I have to admit nothing comes quickly to mind. I did a quick simulation, and the probability (assuming I haven't made any mistakes) comes out to something close to $0.3489$. Doesn't ring any bells yet. I'll keep thinking about it. Interesting problem. – Brian Tung Feb 17 at 5:30
• @BrianTung Thanks and by the way the answer is $3/8$ as stated by my professor. – Rohan Shinde Feb 17 at 5:33
• Weird. I must have done something wrong. Let me make sure I've picked the right PRNG. – Brian Tung Feb 17 at 5:36
• @Digamma: Yeah, it was all roundoff error. I get $3/8$ now, too. – Brian Tung Feb 17 at 5:39
• There is a natural $n$-d generalization (even length sums = odd length sums), and the 4-d probability is 13/48. As in the 3-d case, the measure in two orthants is 1, and the measure in all others is 1/6. My simple symbolic integration routine crashed Mathematica for the 5-d case, however. – Jim Ferry Feb 20 at 5:21
Suppose $$x$$ and $$y$$ are positive and consider the possible values of $$z$$. Because $$|x|+|y|=|x+y|$$ here, we want the equation $$|z| + |x+y+z| = |x+z| + |y+z|$$ to hold.
Assume $$x \le y$$; in this case, we have $$z \le x+z \le y+z \le x+y+z$$, and so we can consider five possibilities based on which of these are positive.
1. $$0 \le z \le x+z \le y+z \le x+y+z$$. Then the equation holds. You already know this case.
2. $$z \le 0 \le x+z \le y+z \le x+y+z$$. Then the equation simplifies to $$x+y=x+y+2z$$, which has probability $$0$$.
3. $$z \le x+z \le 0 \le y+z \le x+y+z$$. Then the equation simplifies to $$x+y=y-x$$, which has probability $$0$$.
4. $$z \le x+z \le y+z \le 0 \le x+y+z$$. Then the equation simplifies to $$x+y = -x-y-2z$$, which has probability $$0$$.
5. $$z \le x+z \le y+z \le x+y+z \le 0$$. Then the equation holds. This case is new.
The same thing happens when $$x \ge y$$, so that doesn't need to be considered separately.
So we see that when $$x$$ and $$y$$ are positive, we want either $$z$$ to be positive as well, or we want $$x+y+z$$ to be negative.
By symmetry, this covers all the possibilities. The equation holds when:
• All three of $$x,y,z$$ are positive;
• Two of $$x,y,z$$ are positive, but $$x+y+z$$ is negative;
• Two of $$x,y,z$$ are negative, but $$x+y+z$$ is positive;
• All three of $$x,y,z$$ are negative.
The regions inside $$[-1,1]^3$$ where these hold have volume respectively:
• $$1$$ (it's a cube of side length $$1$$);
• $$\frac12$$ (it's three pyramids that form a corner of a cube with $$\frac16$$ the volume);
• $$\frac12$$;
• $$1$$.
So total volume $$3$$ (out of $$8$$), so the equation holds with probability $$\frac38$$.
• Can you try to generalize it for n-dimensional space please? – Rohan Shinde Feb 20 at 9:30
• Not with the same approach; the amount of casework grows exponentially. – Misha Lavrov Feb 20 at 14:32
A simulation using R statistical software, for those interested, agrees with the theoretical answer.
> x<-runif(10^7,-1,1)
> y<-runif(10^7,-1,1)
> z<-runif(10^7,-1,1)
> mean(abs(x)+abs(y)+abs(z)+abs(x+y+z)==abs(x+y)+abs(x+z)+abs(y+z))
[1] 0.3749906
The following simulation in Matlab gives simulated probabilities, $$10^7$$ results for n from 2 to 10. Things written after percent signs are comments.
for n=2:10 % dimension n
b=cell(1,n);
[b{:}]=ndgrid(0:1); % n different 2x2x2x...x2 arrays
A=cat(20,b{:});
B=reshape(A,2^n,n); % array contains 0s and 1s to define all 2^n sums
C=(-1).^sum(B,2); % array contains +1 or -1 for each sum
for m=1:10000 % 10000 times 1000 trials is 10^7
D=(2*rand(1000,n)-1)*B'; % 1000 trials, all 2^n sums
E=abs(D)*C; % all 1000 combined sums
k(n)=k(n)+sum(abs(E)<1e-8); % count combined sums =0, give or take underflow
end;
disp(num2str([n k(n)/10000000]));
end;
2.0000 0.4998
3.0000 0.3748
4.0000 0.2710
5.0000 0.1900
6.0000 0.1240
7.0000 0.0815
8.0000 0.0494
9.0000 0.0297
10.0000 0.0168 | 2019-06-19T03:19:43 | {
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https://math.stackexchange.com/questions/3205349/minimum-value-of-4-digit-number-divided-by-sum-of-its-digits | # Minimum value of $4$-digit number divided by sum of its digits?
If x was a positive 4 digit number and you divide it by the sum of its digits to get the smallest value possible, what is the value of x? For example (1234 = 1234/10)
I got 1099 as my answer however I don't know if this is right or how to prove it.
• Computer search shows that if you ask the question for k digit numbers you get the sequence $1,19,199,1099,10999,109999,1099999 ...$ I've verified the pattern out to $k = 8$ (after which I no longer have the patience to run my code). There is a natural conjecture, which could perhaps be proved by induction. – John Coleman Apr 28 '19 at 12:25
• The initial $1$ in the sequence $1, 19, 109, 1099, 10999, ..$ is of course arbitrary since all non-zero 1-digit numbers give the same answer. It isn't too hard to see that the argument that @OldBoy gives below generalizes to all cases which are 4 digits or more. – John Coleman Apr 28 '19 at 14:21
• As a possibly interesting aside, it looks like $1098$ gives the smallest integer quotient when divided by the sum of its digits. – Barry Cipra Apr 29 '19 at 15:37
• @BarryCipra That is an interesting observation (where by integer quotient you only count cases where the sum of the digits divides the number). For 4,5,6,7 digits the solutions are $1098, 10989, 109888, 1078999$. No clear pattern, although they are all less than the corresponding solutions for real-valued quotients. By the way -- I liked your "Misteaks" book and have recommended it to some of my calculus students on more than one occasion. – John Coleman Apr 29 '19 at 16:36
• @JohnColeman, thanks! Your sequence in response to my comment is oeis.org/A034727 at OEIS. – Barry Cipra Apr 29 '19 at 17:58
Let $$x$$ be the 4-digit number ending in digit not equal to 9 with the sum of digits equal to $$s$$. If you increase the last digit by one, the number becomes $$x+1$$ and the sum of digits becomes $$s+1$$.
You can easily prove that:
$$\frac xs>\frac{x+1}{s+1}$$
This is true because $$x>s$$. So to reduce the ratio as much as possible, the last digit has to be as big as possible, which is $$9$$.
Now consider number $$x$$ with the third digit not eqaul to 9. Increase the third digit by one and the number becomes $$x+10$$ and the sum increases by 1. Again, you can show that:
$$\frac xs>\frac{x+10}{s+1}$$
This is true because $$x>10s$$. So the third digit also has to be as big as possibe. So the last two digits must be equal to 9.
Consider now the first digit and decrease it by one. The number becomes $$x-1000$$ and the sum becomes $$s-1$$.
$$\frac xs>\frac{x-1000}{s-1}$$
This is true because it is equivalent to:
$$x<1000s$$
...and you know that $$s$$ is certainly greater than 18 (the last two digits must be 9). So the first digit has to be as small as possible i.e. 1 and the number is of the form:
$$1a99$$
The ratio you want to minimize now becomes:
$$\frac{100a+1099}{a+19}=100-\frac{801}{a+19}$$
Minimum value is reached for minumum value of $$a$$ which is 0. So the solution is: 1099.
The answer 1099 is confirmed by the following MiniZinc model:
% decision variable array
array[1..4] of var 0..9: digits;
var int: objective = sum([digits[i] * pow(10, i-1) | i in 1..4]) div sum(digits);
constraint digits[4] > 0; % otherwise, there would be less than 4 digits
solve minimize objective;
I used Python to prove that shortest posible value is '1099'.
m = []
for eve in range(1000, 10000):
nums = list(map(int, list(str(eve))))
m.append((nums, eve / sum(nums)))
m.sort(key=lambda x: x[1])
print(m[0])
• Golfed: print(min((n/sum(int(d)for d in str(n)),n)for n in range(1000,10000))[1]) 73 bytes :) – Daniel Apr 29 '19 at 4:15
• I believe that Daniel was only aiming to make a joke linked to codegolf.stackexchange.com, since he put a smiley in the end of the comment. I don't think that post aimed to say that the code was better, – Ertxiem - reinstate Monica Apr 29 '19 at 11:33
• @Daniel If you haven't done so already, perhaps you could post this as a question on Code Golf. I suspect that something like J could get this down to less than 10 bytes. – John Coleman Apr 29 '19 at 13:59
Write the $$4$$-digit number as $$x = 1000a+100b+10c+d$$, where $$a,b,c,d$$ are integers from $$0$$ to $$9$$ (and $$a \neq 0$$.)
You want to minimize $$A = \dfrac{1000a+100b+10c+d}{a+b+c+d}$$, which can be written as \begin{align*} A &= \dfrac{999a+99b+9c}{a+b+c+d}+1\\ &\geq \dfrac{999a+99b+9c}{a+b+c+9} + 1 & \text{(since the numerator} > 0\text{, set } d = 9)\\ &= \dfrac{990a+90b-81}{a+b+c+9} + 10\\ &\geq \dfrac{990a+90b-81}{a+b+9+9} + 10 & \text{(since the numerator} > 0\text{, set } c = 9)\\ &= \dfrac{900a-18\times90-81}{a+b+9+9} + 100\\ &= \dfrac{900a-1701}{a+b+18} + 100\\ &\geq \dfrac{900\times 1-1701}{1+b+18} + 100 & \text{(since the } \textbf{denominator} > 0\text{, set } a = 1)\\ &\geq \dfrac{900-1701}{1+0+18} + 100, & \text{(since the numerator} < 0\text{, set } b = 0)\\ \end{align*} and the minimum attains when $$x = 1099$$.
• +1 and this method answers the question for $n$ digit numbers. – Ethan Bolker Apr 28 '19 at 21:59
• Since $18 \cdot 90 + 81 = 1620 + 81 = 1701$, shouldn't you have $\frac{900a - 1701}{a + b + 18} + 100$? – N. F. Taussig Apr 29 '19 at 9:50
• Ahh you’re right. Thanks. – Jerry Chang Apr 29 '19 at 10:36 | 2020-01-20T09:30:17 | {
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https://www.physicsforums.com/threads/finding-the-satellite-orbital-radius.778118/ | # Finding the Satellite Orbital Radius
1. Oct 25, 2014
### logan3
1. The problem statement, all variables and given/known data
If NASA wants to put a satellite in a circular orbit around the sun so it will make 2.0 orbits per year, at what distance (in astronomical units, AU) from the sun should that satellite orbit? The earth's orbit is 1.0 AU from the sun.
$T_E = \frac {1.0orbit}{yr}$
$T_s = \frac {2.0orbit}{yr}$
$r_E = 1.0AU$
2. Relevant equations
$\frac {T_E^2} {r_E^3} = \frac {T_s^2} {r_s^3} => r_s = \sqrt[3]{\frac {T_s^2 r_E^3}{T_E^2}}$
3. The attempt at a solution
Using Kepler’s 3rd law of planetary motion, we know that the orbital period squared is proportional to the radius of the orbit cubed. Since the proportion is approximately the same for all objects orbiting relative to the Sun, then: $\frac {T_E^2} {r_E^3} = \frac {T_s^2} {r_s^3}$
Solving for the radius of the satellite and plugging in the values, we get:
$r_s = \sqrt[3]{\frac {T_s^2 r_E^3}{T_E^2}} = \sqrt[3]{\frac {(\frac {2.0orbit}{yr})^2 (1.0AU)^3}{(\frac {1.0orbit}{yr})^2}} = \sqrt[3]{4.0AU^3} = 1.5874AU \sim 1.6AU$
Thank-you
Last edited: Oct 25, 2014
2. Oct 25, 2014
### ehild
The satellite makes 2 orbits per year, what is its orbital period?
3. Oct 25, 2014
### SteamKing
Staff Emeritus
The further away from the sun, the longer the orbital period. That's Kepler's (unpublished) Fourth Law.
4. Oct 25, 2014
### logan3
$T_s = \frac {2 orbits} {1yr} x \frac {2\pi} {1 orbit} x \frac {1yr} {365days} x \frac {1day} {24hrs} x \frac {1hr} {3600s} = \frac {\pi} {7.884x10^6s}$
$G = 6.67x10^{-11} Nm^2/kg^2$
$M_\odot = 1.9885x10^{30} kg$
$1 AU = 1.4959x10^{11} m$
$T_s^2 = \frac {4\pi^2 r^3}{GM_\odot} \Rightarrow r = \sqrt[3]{\frac {T_s^2 GM_\odot}{4\pi^2}} = 81.102m x \frac {1 AU}{1.4959x10^{11} m} = 5.4216x10^{-10}AU$
5. Oct 25, 2014
### Staff: Mentor
You don't need meters. The approach in post 1 was fine, you just plugged in the wrong orbital period.
And no, an orbital radius of a few meters is certainly wrong (this would be inside the sun!!).
T is the duration of an orbit, not the number of orbits per time (that is the inverse value).
6. Oct 25, 2014
### ehild
2 total orbits last one year, how long time is one orbit? (in years)
7. Oct 25, 2014
### SteamKing
Staff Emeritus
5.4216*10-10AU works out to be about 81 meters. Does this sound like a reasonable answer?
If you make 2 orbits per year, the orbital period is going to be a half year for each orbit.
You know that the earth's orbital period = 1 year (it's kind of a definition) and you have Kepler's third law which relates orbital distance to orbital period.
Try again, and put a little thought into your calculations.
8. Oct 25, 2014
### logan3
For the satellite, one orbit (in years) is half a year. Thus,
$T_E = 1yr$
$T_s = 0.5yr$
$r_E = 1.0AU$
$r_s = \sqrt[3]{\frac {T_s^2 r_E^3}{T_E^2}} = \sqrt[3]{\frac {(0.5 yr)^2 (1.0 AU)^3}{(1 yr)^2}} = \sqrt[3]{0.25 AU^3} = 0.62996 AU \sim 0.63 AU$
In order for the orbital period to decrease, the object must be closer to the Sun. 0.63AU is closer than the Earth and has a shorter orbital period (0.5 yr), therefore, numerically the result corresponds to the aforementioned reasoning.
Using the second method:
$T_s = 0.5 yr x \frac {365 days} {1 yr} x \frac {24 hrs} {1day} x \frac {3600 s} {1 hr} = 1.5768x10^7 s$
$G = 6.67x10^{-11} Nm^2/kg^2$
$M_\odot = 1.9885x10^{30} kg$
$1 AU = 1.4959x10^{11} m$
$T_s^2 = \frac {4\pi^2 r^3}{GM_\odot} \Rightarrow r = \sqrt[3]{\frac {T_s^2 GM_\odot}{4\pi^2}} = 9.4177x10^{11} m x \frac {1 AU}{1.4959x10^{11} m} = 0.62957 AU \sim 0.63 AU$
The answers from both methods match.
Thank-you | 2017-08-17T12:59:19 | {
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https://math.stackexchange.com/questions/966517/how-to-select-the-right-modulus-to-prove-that-there-do-not-exist-integers-a-an | # How to select the right modulus to prove that there do not exist integers $a$ and $b$ such that $a^2+b^2=1234567$?
I understand the solution but I don't know how the author decided to start with modulo 4 instead of something else? What is it about the expression $a^2+b2=1234567$ that would trigger us to select modulo 4 instead of something else.
I tried to solve the question in a similar manner using modulo 2 but eventually got stuck. This leads me to believe that this question can only be solved using modulo 4. Is this true?
• Perhaps it is because the only non-zero quadratic residue modulo $4$ is $1$. Another good candidate would be $8$, and $16$ as well. – Yiyuan Lee Oct 10 '14 at 10:20
• This is just a question of knowing the problem type, I think (sums of squares). My first thought on reading the title was "Try $4$ first, and if that doesn't work, try $8$." – TonyK Oct 10 '14 at 10:22
Considering in mod $4$ is a good 'tool' in such a question just because we have for any integer $a$ $$a^2\equiv 0,1\pmod 4$$ as the author says. This fact is useful in this case just because $$a^2+b^2\not\equiv 3\pmod 4$$ $$1234567\equiv 3\pmod4.$$ Note that this works just because $1234567\equiv 3\pmod 4$.
P.S. Considering in mod $3,8,16$ is also useful.
I'll give you an example. In the following question, considering in mod $3$ may be the first choice (try, and you'll see why) :
Question : Find all positive integers $(n,x,y)$ such that $$y^2=x^n-x+2.$$ The answer is $(n,x,y)=(2m,2,2^m)$ for all positive integers $m$.
Considering in mod $3$ works because $y^2\equiv 0,1\pmod 3$.
• did you choose mod 3 because the expression has 3 unknown terms and could you please direct me to some sites that show how to solve such problems? I've only started reading about modular arithmetic and have never come across such problems – mauna Oct 11 '14 at 3:40
• @mauna: No. If $x\equiv 0,1\pmod 3$, then the RHS$\equiv 2$, which is a contradiction. So, we have $x\equiv 2$, and so on. sorry, but I don't know such sites. How about using google? – mathlove Oct 11 '14 at 10:39
Try mod 2. Doesn't work. Try mod 3. Doesn't work. Try mod 4. It works. Stop.
I would go along with others in the comments that $4$ is the obvious one to choose, since the basic theorems on sums of two squares depend on behaviour mod $4$. The only thing which using $8$ adds is that numbers of the form $8n+7$ cannot be written as the sum of fewer than four squares.
There is a theorem which tells us that $N$ can be written as a sum of two squares only if every prime factor of the form $4n+3$ appears to an even power in the factorisation. So if my first thought had been to try mod $4$, and this had failed, I would be looking to factorise the number into prime factors, or alternatively to spot an easy decomposition into two squares.
I can't see that trying another modulus would make any progress.
HINT:
$$1234567\equiv-1\pmod4\equiv3$$
and for any integer $a,a^2\equiv0,1\pmod4\implies a_1^2+a_2^2\equiv0,1,2\pmod4$
$\quad$ I don't know how the author decided to start with modulo $4$ instead of something else?
$\qquad\qquad\qquad\qquad$
$\qquad\qquad\qquad\qquad\quad$ Hope this picture will clarify all your questions... | 2019-12-10T19:04:52 | {
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https://math.stackexchange.com/questions/3367286/n-identical-objects-in-k-distinct-bins-with-k-1-boxes-having-an-odd-number-of-ob | # n identical objects in k distinct bins with k-1 boxes having an odd number of objects, and one box having x number of objects such that x mod 3 = 2?
How many ways are there to distribute 32 indistinguishable objects into five distinct bins such that four bins have an odd number of objects and one bin has x objects such that x mod 3 = 2.
Initially, I realized that 4 of the boxes must have at least 1 object, and that one box must have at least 2 objects. To maintain the first four objects as odd, we must add pairs of objects to them and to maintain the last bin as producing x mod 3 = 2, we must add triplets of objects to them. But I am unsure where to go from there.
• Well, don't forget that there are odd numbers $\equiv 2\pmod 3$, so one bin might discharge two constraints. – lulu Sep 23 at 21:11
First suppose the question means that bins 1,2,3,4 have an odd number and bin 5 has $$x$$ objects
Your idea regarding initially putting 1 object into each of bins 1 to 4 and 2 into bin 5 is very sensible. Notice also that 4 odd numbers plus $$x$$ makes 32 and so $$x$$ is even i.e. $$x$$ is 2 modulo 6.
You must now put 13 pairs of objects into the bins with either 0, 3, 6, 9 or 12 pairs into bin 5.
It will be easy for you to count the cases now. However, you need to decide precisely what the question is asking:-
@lulu raised the issue of whether one of the odd numbers can also be the $$x$$. Can it?
Also, remember that the bin with $$x$$ objects can be any of bins 1 to 5.
Continuing with the solution. You need to do 5 separate calculations:-
If $$N$$ pairs go into the $$x$$ -box, the number of solutions is the number of ways of distributing 13-$$N$$ identical pairs into 4 distinct bins.
You then multiply the sum of the five numbers by 5 to take account of the position of the $$x$$ box.
We have $$\begin{pmatrix}16 \\ 3\end{pmatrix}+\begin{pmatrix}13 \\ 3\end{pmatrix}+\begin{pmatrix}10 \\ 3\end{pmatrix}+\begin{pmatrix}7 \\ 3\end{pmatrix}+\begin{pmatrix}4 \\ 3\end{pmatrix}$$
which must then be multiplied by 5.
• Any of the boxes is able to be the box full-filling the modulo constraint. The odd numbers can also be x – KengoL Sep 23 at 22:09
• UPDATE: I reread the question, the only constraints is that four boxes are odd and the last box has the 2 mod 3 property. Any distribution of the objects which follows this is valid. – KengoL Sep 23 at 22:18
• When you say "last box" do you mean one other than the odd boxes? – S. Dolan Sep 23 at 22:21
• In that case the method given above works. Are you familiar with "stars and bars" type problems? – S. Dolan Sep 23 at 22:28
• I am familiar. I am just alittle unclear about accounting for the number of pairs which are not 0,3,6,9,12 in box 5. Distributing pairs into each box without taking into the condition would just be (13 + 5 - 1)C(13), but I am still unclear as to how to take into account the # of pairs that cannot be added to the 5th box (of which there are 7) – KengoL Sep 23 at 22:40 | 2019-11-19T02:26:47 | {
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http://mathhelpforum.com/geometry/17924-vectors.html | # Math Help - vectors!!
1. ## vectors!!
hello!!
i'm having a little trouble converting between vector forms!!
so how would you convert the vector V = (3,40 degrees) into cartesian form
and convert v = 2i + 5j into polar form...
thanks!!
2. Hello, sir nerdalot!
How would you convert the vector $v \:= \:(3,\,40^o)$ into cartesian form
and convert $v \:= \:2i + 5j$ into polar form?
If you sketch this diagram, the pieces should come together.
Code:
|
| *
| /|
| r / |
| / |y
| / |
|/θ |
- - * - - + - -
| x
We can see that: . $\begin{array}{c}x \: =\: r\cos\theta \\ y \: =\: r\sin\theta\end{array}$
. . . and that: . $\begin{array}{ccc}r & = & \sqrt{x^2+y^2} \\ \theta & = & \tan^{-1}\!\left(\frac{y}{x}\right)\end{array}$
Then use those conversion formulas . . .
$\begin{Bmatrix}x & = & 3\cos40^o & \approx & 2.3 \\
y & = & 3\sin40^o & \approx & 1.9\end{Bmatrix}\qquad (3,\,40^o)\;\Rightarrow\;2.3i + 1.9j$
$\begin{Bmatrix}r & = & \sqrt{2^2+5^2} & = & \sqrt{29} \\
\theta & = & \tan^{-1}\!\left(\frac{5}{2}\right) & \approx & 68.2^o\end{Bmatrix}\qquad2i + 5j\;\Rightarrow\;\left(\sqrt{29},\,68.2^o\right)
$
3. thanks for that....
but how would i find the angle between these two vectors???
i thought about trying to use the cosine/sine rule but found that there wasn't enough information to do this.
4. Originally Posted by sir nerdalot
thanks for that....
but how would i find the angle between these two vectors???
i thought about trying to use the cosine/sine rule but found that there wasn't enough information to do this.
use the dot product formula
Recall: If $\bold {a,b}$ are two vectors, then $\boxed { \bold {a} \cdot \bold {b} = |\bold {a}||\bold {b}|\cos \theta }$ where $\theta$ is the angle between them.
you have enough information to solve that formula for $\theta$
do you remember how to take the dot product?
5. Hello, sir nerdalot!
I thought about trying to use the cosine/sine rule . . . . Good thinking!
but found that there wasn't enough information to do this.
. .
Well, there is . . . but it takes some Olympic-level gymnastics.
Given two vectors $\vec{u}$ and $\vec{v}$, let $\theta$ be the angle between them.
Code:
Q
*
* \
v * \
* \ w
* \
* θ \
* * * * * * * * * *
O u P
We have: . $\vec{u} \,= \,\overrightarrow{OP},\;\;\vec{v} \,=\, \overrightarrow{OQ},\;\;\theta \,= \,\angle POQ$
. . If $\vec{w} \,= \,\overrightarrow{PQ}$, then: . $\vec{w} \:= \:\vec{u} - \vec{v}$
Using the Law of Cosines, we have: . $|w|^2 \;=\;|u|^2 + |v|^2 - 2|u||v|\cos\theta$ .[1]
A small detour:
. . $|\vec{w}|^2 \;=\;|\vec{u} - \vec{v}|^2 \;=\;(\vec{u}-\vec{v})\cdot(\vec{u}-\vec{v}) \;=\;\vec{u}\cdot\vec{u} - \vec{u}\cdot\vec{v} - \vec{v}\cdot\vec{u} + \vec{v}\cdot\vec{v}$
. . $|w|^2\;=\;|\vec{u}|^2 - 2(\vec{u}\cdot\vec{v}) + |\vec{v}|^2$ .[2]
Substitute [2] into [1]: . $|\vec{u}|^2 - 2(\vec{u}\cdot\vec{v}) + |\vec{v}|^2 \;=\;|\vec{u}|^2 + |\vec{v}|^2 - 2|\vec{u}||\vec{v}|\cos\theta$
This simplifies to: . $2|\vec{u}||\vec{v}|\cos\theta \:=\:2(\vec{u}\cdot\vec{v})\quad\Rightarrow\quad\b oxed{\cos\theta \;=\;\frac{\vec{u}\cdot\vec{v}}{|\vec{u}||\vec{v}| }}$ . . . . There!
6. coolies and so i substitute the two sides into u and v
thanks all!!! | 2014-04-18T06:24:07 | {
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https://math.stackexchange.com/questions/2403367/how-many-6-digit-numbers-contain-exactly-two-even-digits | # How many $6$-digit numbers contain exactly two even digits
How many 6 digit numbers contain exactly two even digits?
I try to solve it by separating it into two cases:
case a: no zero in the beginning (and I assume that the first digit is even number)
-there are $4$ options for the first digit
-for the next digit is $\binom{5}{4}$
-and I assume that in some place there is the next is the next even number
$$4 \binom{5}{4} \cdot 4 \cdot 4 \cdot 4 \cdot 4$$
case b: zero is not allowed to be in the opening digit and I place there an odd number
-there are $4$ options for the odd number (no zero and no even number)
-and for the rest of the digits ($5$ left), I need to place the rest of the numbers so
$$4 \cdot \binom{5}{4} 4 \cdot 4 \cdot 4 \cdot \binom{4}{4}$$
To be honest, I think I really screwed up in this question. Can you help me here?
And explain how to solve this question and this kind of question?
In how many $5$ digit numbers are there digits that appear more than once?
I do not have any idea how to solve this.
• I understand why you have four options for the leading digit if the leading digit is even, but why do you have four options for the other digits? – N. F. Taussig Aug 23 '17 at 11:28
• because there is 4 even numbers 2,4,6,8 am i worng ? – NedyLearing Aug 23 '17 at 11:35
• The number $0$ is even since $0 = 2 \cdot 0$. Also, there are five odd digits, $1$, $3$, $5$, $7$, and $9$. – N. F. Taussig Aug 23 '17 at 11:36
How many six-digit positive integers contain exactly two even digits?
Remember that an integer $n$ is said to be even if there exists an integer $m$ such that $n = 2m$. An integer that is not even is said to be odd. Hence, there are five even digits: $$0, 2, 4, 6, 8$$ and five odd digits: $$1, 3, 5, 7, 9$$
Your strategy of breaking the problem into two cases is correct. Let's correct your calculations for those cases.
Case 1: The leading digit is even.
Since the leading cannot be zero, there are four ways to fill the leading digit. There are five ways to choose the position of the other even digit and five ways to fill that position with an even digit. There are five ways to fill each of the four remaining positions with an odd digit.
$$4 \binom{5}{1} \cdot 5 \cdot 5^4 = 4 \binom{5}{1} 5^5$$
Case 2: The leading digit is odd.
Two of the remaining five positions must be filled with even digits. There are $\binom{5}{2}$ ways to select the positions of the even digits and five ways to fill each of those positions with an even digit. There are five ways to fill each of the four other positions with an odd digit.
$$\binom{5}{2} \cdot 5^2 \cdot 5^4 = \binom{5}{2} 5^6$$
Total: Since the above cases are mutually exclusive and exhaustive, the number of six-digit positive integers that contain exactly two even digits is $$4 \binom{5}{1} 5^5 + \binom{5}{2} 5^6$$
In how many five-digit positive integers are there digits that appear more than once?
Hint: Subtract the number of five-digit positive integers in which all the digits are distinct from the total number of five-digit positive integers.
The total number of five-digit positive integers is $$9 \cdot 10 \cdot 10 \cdot 10 \cdot 10 = 90000$$ since the leading digit can be chosen in nine ways (as $0$ is prohibited) and each of the remaining digits can be selected in $10$ ways. Alternatively, the largest five-digit number is $99999$ and the largest number with fewer than five digits is $9999$, so the number of five-digit positive integers is $$99999 - 9999 = 90000$$ The number of five-digit positive integers with distinct digits is $$9 \cdot 9 \cdot 8 \cdot 7 \cdot 6 = 27216$$ since we have nine choices for the leading digit (as $0$ is prohibited), nine choices for the thousands digit (since we cannot use the leading digit), eight choices for the hundreds digit (since we cannot use the tens thousands or thousands digits), seven choices for the tens digit, and six choices for the units digit. Therefore, the number of five-digit positive integers that do contain a digit that appears more than once is $$9 \cdot 10 \cdot 10 \cdot 10 \cdot 10 - 9 \cdot 9 \cdot 8 \cdot 7 \cdot 6 = 90000 - 27216 = 62784$$
• thank you very much :) i understund it finaly! :) for the second one i will try it right away :) – NedyLearing Aug 23 '17 at 15:34
• for the second quation is my way of getting the answer is correct ? all the posibalits (10^5) to get five digit numbers - posibalits with no repeation at all = my answer ? – NedyLearing Aug 23 '17 at 16:42
• You did not take into account the fact that the leading digit cannot be zero. I have added a solution to the second problem. You can read it by hovering over the shadow box. – N. F. Taussig Aug 23 '17 at 16:59
• oh god .. i allways do mistakes like that .. =\ but i think now i get it i must say realy thank you for all your help youre awsome ! – NedyLearing Aug 23 '17 at 17:06 | 2019-09-15T22:34:14 | {
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https://math.stackexchange.com/questions/2402451/what-exactly-do-we-call-a-linear-equation | # What exactly do we call a linear equation?
I know there's a difference between the linear algebra "linear function" and the "affine function", but it confused me about what we consider a "linear equation" -precisely how many variables allowed, what order could a variable carry?
• A linear equation is an equation of the form $\sum_{k=1}^n a_kx_k = c$ for variables $\{x_k\}_{k=1}^n$ and constants $\{a_k\}_{k=1}^n$ and $c$. – Michael L. Aug 22 '17 at 15:05
• Thanks a lot, what does "c" means in the sum? is it a parameter? – lun Aug 22 '17 at 16:48
• As I said, it is an arbitrary constant. – Michael L. Aug 22 '17 at 20:33
• I don't know what "what order could a variable carry" means. – Gerry Myerson Aug 27 '17 at 2:54
Here are various comments on linear functions versus linear equations.
First there is the idea of a linear combination. Essentially, linear combination means a "sum of multiples" of something. For example
$$3x-2y+4z$$
is a linear combination of the variables $x,y$ and $z$, that is, a sum of multiples of $x,y$ and $z$.
If you set this expression equal to a constant value such as
$$3x-2y+4z=7$$
then you have a linear equation.
A linear function, on the other hand, is a function with the property that the function of a linear combination of two or more variables equals the same linear combination of the function of the two variables. This definition is difficult to understand without an example.
Let $f(x)=3x$. Then find $f(4a+7b)$
Solution:
$$f(4a+7b)=3(4a+7b)=12a+21b$$
Now look at the same linear combination of $f(a)$ and $f(b)$.
$$4f(a)+7f(b)=4(3a)+7(3b)=12a+21b$$
It is the same. That is
$$f(4a+7b)=4f(a)+7f(b)$$
This notion of linearity occurs in many areas of mathematics.
Notice, for example, that definite integrals are linear. For example
$$\int_0^15x-6x^2\,dx=5\int_0^1x\,dx-6\int_0^1x^2\,dx$$
Differentiation is also linear. For example
$$\frac{d}{dx}\left(3\sin x+5e^x\right)=3\frac{d}{dx}\sin x+5\frac{d}{dx}e^x$$
Should, at some point, you study Laplace transforms you will find that they are linear transforms.
$$\mathcal{L}\{a\cdot f(t)+b\cdot g(t)\}=a\cdot\mathcal{L}\{f(t)\}+b\cdot\mathcal{L}\{g(t)\}$$
• In fact a linear equation is nothing but $f(x)=y$ for some linear function $f$. – user856 Aug 27 '17 at 1:53
• @Rahul I was afraid that I had already gone on at too great a length. – John Wayland Bales Aug 27 '17 at 1:58 | 2021-01-28T11:29:29 | {
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https://math.stackexchange.com/questions/3993012/connecting-the-two-definitions-of-e/3993026 | # Connecting the two definitions of $e$
There are two different ways I've been taught to understand the meaning of $$e$$.
1. $$e$$ can be defined as the total growth from continuously compounding interest in a single period. To make the growth rate the same as the dollar amount, consider I start with $$\1$$. Then, $$a_n$$ represents the amount I have after compounding $$n$$ times. \begin{align} a_n&=1\cdot\left(1+\frac{1}{n}\right)^n\\[5pt] e=a_{\max}&=\lim_{n \to \infty}\left(1+\frac{1}{n}\right)^n \end{align} As we increase the number of times we compound the interest, the final amount (in the case that we start with $$\1$$), approaches $$e$$.
2. All exponential functions increase at a rate proportional to their current value. For example, consider the derivative of $$2^x$$: \begin{align} \frac{d\left(2^x\right)}{dx}&=\lim_{h \to 0}\frac{2^{x+h}-2^x}{h}\\[5pt] &=\lim_{h \to 0}\frac{2^x2^h-2^x}{h}\\[5pt] &=\lim_{h \to 0}\frac{2^x\left(2^h-1\right)}{h}\\[5pt] &=2^x\lim_{h \to 0}\frac{\left(2^h-1\right)}{h} \end{align} From this we can see that the rate at which $$2^x$$ changes is proportional to its current value, with a proportionality constant of $$\lim_{h \to 0}\frac{\left(2^h-1\right)}{h}.$$ For $$2^x$$, this constant happens to be $$\ln(2)$$ (why is this the case?). $$e$$ can be defined as the base where this proportionality constant is $$1$$; therefore, the rate $$e^x$$ grows at is exactly its current value.
My question is this: how can the same constant have two (apparently) separate definitions? Is there a way of looking at $$e$$ that will combine these two definitions, and make it obvious as to why they are true? Also, why is the proportionality constant of any exponential functions growth rate the natural log of its base? Is there a way to understand why this is the case intuitively?
• $\exp(1)=\sum\limits_{n=0}^\infty \dfrac{1}{n!}$ is another possible definition – Henry Jan 20 at 16:04
• @Henry Yes, this makes sense as it's the Taylor Series expansion of $e$ – John Hippisley Jan 20 at 16:09
• "how can the same constant have two (apparently) separate definitions?" Happens all the time. Definition 1: $k$ is the smallest even prime. Definition 2: $m$ is the sum of the multiplicative identity with itself. Definition 3: $j$ is the square root of the smallest composite natural number. We can calculate that $k= 2$ and $m=2$ and $j=2$. And we could analysis and figure out why those must be equal... or we could accept they are. – fleablood Jan 20 at 16:36
• From an intuitive perspective, the two characterisations of the exponential function should definitely be related, if not equivalent. The second characterisation of $e^x$ that you gave shows that the exponential growth of a quantity is proportional to the quantity itself. This is the same as in the compound interest example, where the more money you have in the bank, the faster the rate at which it grows. Then taking things to the limit means that you are dealing with continuous (as opposed to discrete) growth. – Joe Jan 22 at 16:06
Let's start with your second definition. We want to find the derivative of $$f(x) = a^x$$, for some $$a>0.$$ Using the definition of derivative, we get
\begin{align} f'(x) &= \lim_{h\to 0}\frac{f(x+h) - f(x)}{h}\\ &= \lim_{h\to 0}\frac{a^{x+h} - a^{x}}{h}\\ &=a^{x}\lim_{h\to 0}\frac{a^{h} - 1}{h}. \end{align}
Now, as you said, one way to define $$e$$ is to define it such that $$\lim_{h\to 0}\frac{e^{h} - 1}{h} = 1.\tag{*}$$ Now, I'm going to be somewhat informal here. Suppose we have $$\frac{e^{h} - 1}{h} = 1.$$ Then, doing some manipulation: \begin{align} \frac{e^{h} - 1}{h} &= 1\\ e^{h} - 1 &= h\\ e^{h} &= h + 1\\ e &= (h+1)^{1/h} \end{align}
Now, by taking limits we have that $$\begin{equation*} \lim_{h\to 0}\frac{e^{h} - 1}{h} = 1 \iff e = \lim_{h\to 0}(h+1)^{1/h}. \end{equation*}$$
Taking the new limit, $$e = \lim_{h\to 0}(h+1)^{1/h}$$, and setting $$n = 1/h$$, we note that $$h\to 0 \implies n\to \infty$$, and we can change variables to get $$e = \lim_{n\to\infty}\left(1 + \frac{1}{n}\right)^{n},$$ which matches your first definition.
As to understanding where the natural log comes from, i.e. why we have $$\lim_{h\to 0}\frac{a^{h}-1}{h} = \ln(a),$$ this is a straightforward consequence of $$(*)$$: \begin{align} \lim_{h\to 0}\frac{a^{h} - 1}{h} &= \lim_{h\to 0}\frac{e^{\ln(a^{h})}-1}{h}\\ &=\lim_{h\to 0}\frac{e^{h\ln(a)}-1}{h}\\ &=\frac{\ln(a)}{\ln(a)}\lim_{h\to 0}\frac{e^{h\ln(a)}-1}{h}\\ &=\ln(a)\lim_{h\to 0}\frac{e^{h\ln(a)}-1}{h\ln(a)}.\\ \end{align} Now, set $$k = h\ln(a)$$, and note that because $$a > 0$$, we have that $$h\to 0\implies k \to 0$$ as well, and so we have by $$(*)$$ that: \begin{align} \lim_{h\to 0}\frac{a^{h} - 1}{h} &=\ln(a)\lim_{h\to 0}\frac{e^{h\ln(a)}-1}{h\ln(a)}\\ &=\ln(a)\lim_{k\to 0}\frac{e^{k} - 1}{k}\\ &=\ln(a)\cdot 1\\ &=\ln(a). \end{align}
This Wikipedia article: Characterizations of the exponential function, shows six ways to define the exponential function $$e^x$$ (thus for the constant $$e$$) and a detailed discussion of equivalency of all the six definitions.
This happens a lot in mathematics. The same mathematical object may have many properties that can characterize the object. Once they are proved to be equivalent to the definition, they can all be regarded as a definition. More generally, there are equivalent definitions of mathematical structures.
For an intuitive understanding of the constant $$e$$ and exponential growth, you may read this article https://betterexplained.com/articles/an-intuitive-guide-to-exponential-functions-e/ from the "BetterExplained" website.
Take the second definition:
Consider: $$\lim_{h\to 0}\frac {b^h -1}h = \ln b = 1$$. And solve for $$b$$. That should be $$e$$ and we could define, supposedly, $$e$$ as the $$b$$ where that limit is $$1$$.
....
$$\lim_{h\to 0}\frac {b^h -1}h =$$
$$\lim_{\frac 1n\to 0}\frac {b^{\frac 1n} -1}{\frac 1n} =$$
$$\lim_{n\to \infty} n(b^{\frac 1n} - 1)=1$$
Now for any $$k$$ if $$k(b_k^{\frac 1k} - 1)=1$$ then $$b_k= (1+ \frac 1k)^k$$ and so if we wave our hands really fast and point out that "Look, it's Halley's comet" it should be reasonable that
If $$\lim_{n\to \infty}n(b^{\frac 1n} -1) = 1$$ then
$$b = \lim_{n \to \infty}(1 + \frac 1n)^n$$.
So the definitions are consistent.
Here is an alternative approach: if we first define the natural logarithm as $$\log x = \int_{1}^{x}\frac{1}{t} \, dt$$ and the exponential function as the inverse of $$\log$$, then it is not much work to show that $$e^x$$ is its own derivative. (I go into more detail about this here.) Then, the limit $$\lim_{n \to \infty}\left(1+\frac{1}{n}\right)^n$$ can be rewritten as $$\lim_{n \to \infty}e^{n\log(1+1/n)} \, .$$ If we make a change of variables $$h=1/n$$, then this limit is equal to $$\lim_{h \to 0^+}e^{1/h\log(1+h)}$$ For small $$h$$, $$\log(1+h) \approx h$$. This can be easily justified by considering the graph of $$\log$$:
$$\log x$$ denotes the area of the region bounded by the hyperbola $$y=1/t$$, the $$t$$-axis, and the vertical lines $$t=1$$ and $$t=x$$. When $$x = 1+h$$ (and $$h$$ is small), this region is approximately rectangular, with a width of $$h$$ and a height of $$1$$. Hence, $$\log(1+h) \approx h$$, with $$h$$ being a slight overestimate of the region's area.
This provides us with a nice heuristic argument that $$\lim_{h \to 0^+}e^{1/h\log(1+h)} = e^{1/h \cdot h} = e \, .$$ | 2021-02-28T10:56:56 | {
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https://math.stackexchange.com/questions/1957336/proving-limit-implies-another-limit | # Proving limit implies another limit
I'm working out some study problems that I don't have answers to yet. Could someone help me out with this one please?
Let $a_n$ and $b_n$ be two positive sequences such that $$b_n = \frac{a_1+a_2...+a_n}{n}.$$ Prove that $\lim_{n\to \infty} a_n=0$ implies $\lim_{n\to \infty} b_n=0$. Is the converse true?
Given $\epsilon>0$, we know there is some $N$ such that $a_n < \epsilon/2$ for all $n > N$. Having chosen this $N$, there exists another $N'\ge N$ such that $\frac{a_1+\cdots+a_N}{n} < \epsilon/2$ for all $n \ge N'$. Then for $n>N'$, $$b_n = \frac{a_1+\cdots+a_N}{n} + \frac{a_{N+1}+\cdots+a_n}{n} < \epsilon/2 + \frac{n\epsilon/2}{n} = \epsilon.$$
Converse: consider the sequence $a_n=(-1)^n$. Then $b_n = \begin{cases}0 & \text{$n$even}\\ -1/n & \text{$n$odd}\end{cases}$ and tends to zero, while $a_n$ does not converge.
To prove the desired convergence, you can use the Cauchy criterion; i.e., use the fact that $a_n$ is Cauchy to examine the differences $$b_{m} - b_{n}.$$
The converse is false: take $a_{n} = (-1)^n$. | 2019-05-21T02:38:13 | {
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http://math.stackexchange.com/questions/371507/calculating-recursive-sequence-numbers-using-generating-functions | # Calculating recursive sequence numbers using generating functions
I have a sequence created with the following recursive formula:
$a_{n+1} = 2 a_{n} + n$
I need to find a generating function for this one, so that I get the formula from which I can calculate the numbers. So according to a book "Generatingfunctionology" reccomended by one user here, I need to do the following steps:
1) Multiply the formula by $x^{n}$, so I get:
$x^{n}a_{n+1} = 2 x^{n}a_{n} + n x^{n}$
2) "Sum over all the values of n for which the relation is valid" - so basically I'm getting (let me omit some steps):
$\sum_{n=0}^{ \infty } a_{n} x^{n} = \frac{1 - 2x + 2x^{2}}{(1-x)^{2}(1-2x)}$
My question is - does this pattern applies to all sequences given by a recursive formula?
And one more thing - this is the result of summing up infinitive sumber of single coefficients, or however you call this part right after the Epsilon sign. How am I supposed to get the sequence elements from this thing?
Or maybe I'm messing everything up?
-
Oh crap, I don't how how to apply latex to my post :/ – khernik Apr 24 '13 at 14:01
Replace your $$ and $$ with $ for in-line tex, or $$ for tex on a new line. For instance \sqrt{x} gives \sqrt{x} and $$\sqrt{x}$$ gives$$\sqrt{x}$$– Clive Newstead Apr 24 '13 at 14:06 Thank you, edited. – khernik Apr 24 '13 at 14:07 does this pattern applies to all sequences given by a recursive formula?... Unsurprisingly, no, but I seem to remember that the book you cite spends quite some time explaining which ones are OK. – Did Apr 24 '13 at 15:19 ## 1 Answer Once you have the generating function, usually the right strategy is to use partial fractions. In this case, you have$\frac{2x^2-2x+1}{(1-2x)(1-x)^2}=\frac{-1}{(1-x)^2}+\frac{2}{1-2x}$. These are now, hopefully, familiar generating functions. We have$\frac{1}{(1-x)^2}=\sum (n+1)x^n$and$\frac{1}{1-2x}=\sum 2^nx^n$. Combining with the partial fraction representation, we get$a_n=-(n+1)+2\cdot 2^n$. Note that the coefficients$-1, 2\$ are from the partial fractions.
In answer to your question of whether this will always work, unfortunately sometimes the generating function doesn't look like anything familiar. However this method does work for many types of recurrences.
-
Thanks! Is there any list of these familiar generating functions? – khernik Apr 24 '13 at 14:23
Look e.g. at Wilf's "generatingfunctionology" (an old edition is free on the ǹet) – vonbrand Apr 24 '13 at 15:32
Also, a few basic ones are in the Wikipedia article, near the end of section 1.1. – vadim123 Apr 24 '13 at 17:38 | 2015-05-29T17:13:55 | {
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https://www.physicsforums.com/threads/torque-and-tension-question.127157/ | Torque and tension question
1. Jul 25, 2006
kankerfist
The picture I attached to this thread is a system that I came across while studying for my final exam today. It is a system that is in equilibrium, and the only numerical values are:
Angle of wire/ground: 30
Angle of strut/ground: 45
Mass of hanging weight: 500 lb
Mass of strut: 100 lb (uniform mass)
I have written out all the forces that I can find, several of which are couples that cancel out, but I cannot find a numerical value for the tension in the cable T. The question asks what the tension is, and I have begun to think there is no way to describe tension numerically from the given data. Can someone please help me out by confirming that there isn't enough data here to describe tension numerically? Thanks a lot!
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2. Jul 25, 2006
Staff: Mentor
Seems like there's enough information to me. Have you applied the equilibrium conditions to the strut? (What's the net torque on the strut due to all the forces on it?)
3. Jul 25, 2006
kankerfist
The net torque at P has to be zero. I guess my problem is that I haven't figured out how to numerically describe the positions of the forces relative to P. For instance, one of the torque forces is the 500 lb weight. I know the direction of the force is straight down, but I'm not sure how to numerically find its position vector relative to P. If I did, then I could calculate the cross product to find the torque due to the 500 lb weight. Thanks for the reply!
4. Jul 25, 2006
Staff: Mentor
Perhaps you are thinking that you need the length of the strut? If that's the problem, just call the length "L". You'll find that you don't need the numerical value: when you write the torque equation, the length will cancel out.
(Tip: Try to solve the problem symbolically. Wait until the last step before reaching for the calculator.)
5. Jul 25, 2006
PPonte
The torque force of the 500 lb weight is applied at the top of the strut and the direction of this force is as you said straight down. Then, you need to find the lever arm.
The lever arm is the perpendicular distance from the axis of rotation to a line drawn along the direction of the force.
Once you know the lever-arm, you may find the torque.
You may follow the same strategy to find the other torques.
Remember Doc Al's hints!
6. Jul 25, 2006
kankerfist
You guys are about 10x as helpful as my teacher! Here is where I have gotten:
Grav Torques from the masses + Torque from wire's tension = 0,
so:
x_strut is distance from P to the 500lb mass lever arm
x_wire is distance from where the wire is grounded to P
h is the height of where the wire and strut meet
(226.8)(9.8)(x_strut) + (45.36)(9.8)(x_strut / 2) = (Tension)(x_wire)
Then geometrically I know that:
Sin[30] = h/(x_strut + x_wire)
Sin[45] = h/x_strut
so:
h = (x_strut + x_wire)Sin[30] = (x_strut)Sin[45]
Solving for x_wire:
x_wire = 0.414 x_strut
Then plugging into the original torque equation gives me:
(226.8)(9.8)(x_strut) + (45.36)(9.8)(x_strut / 2) = (Tension)(0.414 x_strut)
Solving this for Tension yields 5905 N of tension in the wire. Any hints on whether or not I'm approaching this correctly would be greatly appreciated!
Last edited: Jul 26, 2006
7. Jul 26, 2006
kankerfist
Actually, I got the lever arm on tension messed up there. Looks like the lever arm of tension is actually (0.414)(x_strut)Sin(30)
This fix yields a tension of 11816 N of tension in the wire
8. Jul 26, 2006
To provide a clearer definition, it would be better to say that $$x_{strut}$$ is the HORIZONTAL distance from P to the vertical wire connecting the 500lb mass. Also, h is the height of the point where the wire and strut meet.
Really? Check again!
Last edited: Jul 26, 2006
9. Jul 26, 2006
Staff: Mentor
Looks like you have things a bit messed up. Let's do it step by step. (I'll use "L" as the length of the strut.) First, let's find the clockwise torques about point P:
-due to the weight of the strut: (L/2)*(100 lbs)* sin(45)
-due to the hanging mass: (L)*(500 lbs)* sin(45)
Now find the counterclockwise torque about point P, which is due to the tension in the wire. To do that you need to figure out the angle that the wire makes with the strut. It's easy--just use what you know about triangles. I'll call that angle "theta":
-due to wire: (L)*(T) sin(theta)
Once you figure out theta, just set total torque equal to zero (or clockwise torques equal to counterclockwise torques -- same thing) and solve for T. | 2016-12-05T16:40:52 | {
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https://math.stackexchange.com/questions/2499083/understanding-the-c-d-f-of-y-n-max-x-1-x-n-and-y-1-min-x-1-x-n | # Understanding the c.d.f of $Y_n=\max\{X_1,..,X_n\}$ and $Y_1=\min\{X_1,..,X_n\}$
Suppose that $X_1 , . . . , X_n$ form a random sample of size $n$ from the uniform distribution on the interval [0, 1]. Define a new random variable $Y_n=\max\{X_1,..,X_n\}$.
I know that i can compute the c.d.f. of $Y_n$ as follows $$P(Y_n<y)=P(\max\{X_1,..,X_n\}<y)=\prod_{j=1}^n P(X_j<y)= [P(X_j<y)]^n$$
I saw in my book that the c.d.f. of $Y_1$ can be written as $P(Y_1<y)=[1-P(X_j<y)]^n$. But why? I know that the $\max\{\}$ and $\min\{\}$ applications are kind of inverse one each other. But in my mind $Y_n$ and $Y_1$ have the same distribution.
Why $P(Y_1<y)$ is not like this $P(Y_1<y)=P(\min\{X_1,..,X_n\}<y)=\prod_{j=1}^n P(X_j<y)$ ?
I don't realize the consequence of having $\max\{\}$ or $\min\{\}$.
• Yes. I did it because the $X_k$'s are taken at random. – Gabriel Sandoval Nov 1 '17 at 2:10
• As for why $Y_{(n)}$ and $Y_{(1)}$ don't have the same distribution: Consider, as an example, drawing $100$ copies of a uniform random variable on $[0, 1]$. Note that $Y_{(n)}$ should be rather close to $1$ and $Y_{(1)}$ should be rather close to $0$. In particular, it is exceedingly unlikely that $Y_{(n)}$ will be less than $0.5$ or that $Y_{(1)}$ will be more than $0.5$. – Aaron Montgomery Nov 1 '17 at 2:44
If $X_1, X_2, \ldots, X_n$ is your sample, and $Y_1 = \min\{X_1, X_2, \ldots, X_n\}$, then $$\Pr[Y_1 \le y] \ne \prod_{j=1}^n \Pr[X_j \le y],$$ because the right hand side requires all of the $X_j$ to be less than or equal to $y$, whereas the left hand side only requires the smallest of the $X_j$ to be less than or equal to $y$. Thus the left hand side counts events that the right hand side does not. For example, if $n = 5$ and we have $(X_1, X_2, X_3, X_4, X_5) = (0.5, 0.3, 0.2, 0.8, 0.4)$, then $Y_1 = 0.2$. But the event $Y_1 \le 0.6$ is true, whereas the event $X_4 \le 0.6$ is false.
So this might make you ask why such an argument does work for $Y_n = \max\{X_1, X_2, \ldots, X_n\}$? Well, the event $Y_n \le y$ requires the largest of the $X_j$ to be less than or equal to $y$, which in turn implies that all of the $X_j$ are less than or equal to $y$, owing to the fact that $Y_n$ is by construction the largest member of the sample. This is in contrast to the same reasoning applied to $Y_1$: simply saying that the smallest member of the sample is less than or equal to some given number $y$, does not guarantee that the other members of that sample are also less than or equal to $y$. And that's the difference.
How would we modify the successful reasoning for computing $\Pr[Y_n \le y]$ to compute $\Pr[Y_1 \le y]$? The answer is to understand that the natural ordering of the minimum member as being the smallest member of the sample--thus every other member is at least as large--means we should consider the complementary probabilities: $$\Pr[Y_1 > y] = \prod_{j=1}^n \Pr[X_j > y].$$ And if we want to write things out in terms of cumulative distribution functions, we would just say $$\Pr[Y_1 \le y] = 1 - \Pr[Y_1 > y] = 1 - \prod_{j=1}^n \Pr[X_j > y] = 1 - \prod_{j=1}^n (1 - \Pr[X_j \le y]).$$
Having $\max\{X_1, \dots, X_n\} \leq y$ is equivalent to having $X_1 \leq y$ and $X_2 \leq y$ and $\dots$ and $X_n \leq y$. This is why it's convenient to work with in terms of distributions; the fact that we can express this as several independent events separated by the word "and" lets us multiply. But, the same can't be said for the minimum.
What we could say for the minimum is: having $\min\{X_1, \dots, X_n\} \leq y$ is equivalent to have $X_1 \leq y$ or $X_2 \leq y$ or $\dots$ or $X_n \leq y$. However, the word "or" is not nearly as convenient as the word "and" when it comes to probabilities. So, we use a trick instead; we observe that having $\min\{X_1, \dots, X_n\} \geq y$ is equivalent to having $X_1 \geq y$ and $X_2 \geq y$ and $\dots$ and $X_n \geq y$. This is why there are $(1-)$ terms both inside and outside the CDF expressions you presented; we're dealing with the opposite thing typically encoded by a CDF.
• As Jack noted in the comments to the original post, I implicitly assumed that your random variables were independent. Based on the formulas you presented, that should be the case. – Aaron Montgomery Nov 1 '17 at 2:05
The event for the maximum of a sequence of random variables being less than or equal to a value, is the event for all of the variables to be so. When the variables are independent, the probability for the conjoint event is the product of the probabilities of the events. When the events are identically distributed, the probabilities are identical.
\begin{align}\mathsf P(Y_n\leqslant y) {}&=\mathsf P(\max{\{X_i\}}_{i=1}^n\leqslant y) \\[1ex]&= \mathsf P(\bigcap_{i=1}^n\{X_i\leqslant y\}) \\[1ex] &= \prod_{i=1}^n\mathsf P\{X_i\leqslant k\} \\[2ex] \mathsf P\{Y_n\leqslant y\} {}&= \mathsf P\{X_1\leqslant k\}^n\\[1ex] &= y^n~\mathbf 1_{y\in[0;1]}\end{align}
The event for the minimum of a sequence of random variables being greater than a value, is the event for all of the variables to be so. Everything else follows as above, except that we are dealing with the complements of the events whose probability we seek and know.
\begin{align}\mathsf P\{Y_1\leqslant y\} {}&=\mathsf P(\min{\{X_i\}}_{i=1}^n\leqslant y) \\[1ex]&=1 -\mathsf P(\min\{X_i\}_{i=1}^n> y) \\[1ex] & = 1-\mathsf P(\bigcap_{i=1}^n\{X_i>y\}) \\[1ex] & = 1-\prod_{i=1}^n\mathsf P\{X_i>y\} \\[1ex] &= 1-\prod_{i=1}^n(1-\mathsf P\{X_i\leqslant y\}) \\[2ex]\mathsf P\{Y_1\leqslant y\} {} & = 1-{\big(1-\mathsf P(X_1\leqslant y)\big)}^n\\[1ex] &= (1-(1-y)^n)~\mathbf 1_{y\in[0;1]}\end{align}
Notice: This is not quite what you claim the book says, though it is likely what the book actually said.
But in my mind $Y_n$ and $Y_1$ have the same distribution.
Clearly the value for $Y_n$ will certainly be at least as great as that for $Y_1$. Since $Y_n\geqslant Y_1$ almost surely, therefore they cannot follow the same distribution. The least value in the sample is more likely than the greatest, to be smaller than some contant. $$\mathsf P(Y_1{\leqslant}y)\geq\mathsf P(Y_n{\leqslant} y)$$ | 2020-02-25T13:21:49 | {
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Over the past 7 weeks, the Smith family had weekly grocery bills of $74,$69, $64,$79, $64,$84, and $77. What was the Smiths' average (arithmetic mean) weekly grocery bill over the 7-week period? A.$64
B. $70 C.$73
D. $74 E.$85
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Over the past 7 weeks, the Smith family had weekly grocery bills of $7 [#permalink] ### Show Tags 27 Jun 2018, 17:28 7 11 Bunuel wrote: Over the past 7 weeks, the Smith family had weekly grocery bills of$74, $69,$64 $79,$64, $84, and$77. What was the Smiths' average (arithmetic mean) weekly grocery bill over the 7-week period?
A. $64 B.$70
C. $73 D.$74
### Show Tags
29 Jun 2018, 05:58
29
9
First, choose a nice, round number, such as 70, within the range of values.
Then calculate the average with the following formula:
Average = Nice number + Average of differences from the nice number
Average = 70 + (4 - 1 - 6 + 9 - 6 + 14 + 7)/7 = 70 + 21/7 =73
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Re: Over the past 7 weeks, the Smith family had weekly grocery bills of $7 [#permalink] ### Show Tags 28 Jun 2018, 17:07 1 Bunuel wrote: Over the past 7 weeks, the Smith family had weekly grocery bills of$74, $69,$64, $79,$64, $84, and$77. What was the Smiths' average (arithmetic mean) weekly grocery bill over the 7-week period?
A. $64 B.$70
C. $73 D.$74
E. $85 We can determine the average using the formula: average = sum / number: average = (74 + 69 + 64 + 79 + 64 + 84 + 77)/7 = 511/7 = 73 Answer: C _________________ # Scott Woodbury-Stewart Founder and CEO [email protected] 181 Reviews 5-star rated online GMAT quant self study course See why Target Test Prep is the top rated GMAT quant course on GMAT Club. Read Our Reviews If you find one of my posts helpful, please take a moment to click on the "Kudos" button. VP Joined: 14 Feb 2017 Posts: 1369 Location: Australia Concentration: Technology, Strategy GMAT 1: 560 Q41 V26 GMAT 2: 550 Q43 V23 GMAT 3: 650 Q47 V33 GMAT 4: 650 Q44 V36 GMAT 5: 650 Q48 V31 GMAT 6: 600 Q38 V35 GMAT 7: 710 Q47 V41 GPA: 3 WE: Management Consulting (Consulting) Re: Over the past 7 weeks, the Smith family had weekly grocery bills of$7 [#permalink]
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19 Nov 2018, 15:13
1
Admittedly I made a stupid mistake and I organised the numbers in ascending order then selected the median (74) instead of solving Sum/# Terms
The method I used is only for consecutive integers and the punishment answer D is there for suckers like me.
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Re: Over the past 7 weeks, the Smith family had weekly grocery bills of $7 [#permalink] ### Show Tags 25 Jun 2018, 04:10 Solution Given: • Over the past 7 weeks, the Smith family had weekly grocery bills of$74, $69,$64, $79,$64, $84, and$77
To find:
• Smith’s average weekly grocery bill over the 7-week period
Approach and Working:
• Smith’s total bill amount over the 7-week period = (74 + 69 + 64 + 79 + 64 + 84 + 77) = 511
• Therefore, Smith’s average bill amount over the same period = $$\frac{511}{7}$$ = 73
Hence, the correct answer is option C.
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28 Aug 2018, 02:42
Bunuel wrote:
Over the past 7 weeks, the Smith family had weekly grocery bills of $74,$69, $64,$79, $64,$84, and $77. What was the Smiths' average (arithmetic mean) weekly grocery bill over the 7-week period? A.$64
B. $70 C.$73
D. $74 E.$85
NEW question from GMAT® Official Guide 2019
(PS07369)
Easy way out is consider avg 70 -
so diff for all terms is 4, -1 , -6, 9 , -6,14, 7 = sum is 21 = divided by 7 = equal to 3 so 70 +3 = 73 avg
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Re: Over the past 7 weeks, the Smith family had weekly grocery bills of $7 [#permalink] ### Show Tags 10 Sep 2018, 06:34 Bunuel wrote: Over the past 7 weeks, the Smith family had weekly grocery bills of$74, $69,$64, $79,$64, $84, and$77. What was the Smiths' average (arithmetic mean) weekly grocery bill over the 7-week period?
A. $64 B.$70
C. $73 D.$74
E. $85 NEW question from GMAT® Official Guide 2019 (PS07369) $$\frac{( 70 + 4 ) + ( 70 - 1 ) + ( 60 + 4 ) + ( 80 - 1 ) + ( 60 + 4 ) + ( 80 + 4 ) + ( 70 + 7 )}{7}$$ = $$\frac{( 70 + 70 + 60 + 80 + 60 + 80 + 70 ) + ( 4 - 1 + 4 - 1 + 4 + 4 + 7 )}{7}$$ = $$\frac{490 + 21}{7}$$ = $$\frac{490}{7} + \frac{21}{7}$$ = $$70 + 3$$ = $$73$$, Answer must be (C) _________________ Thanks and Regards Abhishek.... PLEASE FOLLOW THE RULES FOR POSTING IN QA AND VA FORUM AND USE SEARCH FUNCTION BEFORE POSTING NEW QUESTIONS How to use Search Function in GMAT Club | Rules for Posting in QA forum | Writing Mathematical Formulas |Rules for Posting in VA forum | Request Expert's Reply ( VA Forum Only ) Intern Joined: 02 May 2018 Posts: 12 GMAT 1: 620 Q46 V29 Re: Over the past 7 weeks, the Smith family had weekly grocery bills of$7 [#permalink]
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10 Sep 2018, 06:46
We can rule out
64, 74 and 85 as they are the extreme numbers.
That leaves just two choices to choose from 70 and 73.
Only 3 values in the 60s and therefore less likely for 70.
73 is the most likely choice. Confirm 73 by taking difference and summing it up
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27 Jan 2020, 11:45
Notice that all options' values have different unit digit;
To Solve: Add the unit digits of each values(4+9+4+9+4+4+7) = 41 and check the unit digit of output (unit digit 1).
for division by 7, 3 should be unit digit of the quotient.
Only Option C satisfies, it is $73! _________________ Unable to Give Up! Re: Over the past 7 weeks, the Smith family had weekly grocery bills of$7 [#permalink] 27 Jan 2020, 11:45
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Join Association For Psychological Science, Fishing Streamers Upstream, Iberia Bank Checking Account, It's Not The End Of The World Video, Bat 21 Full Movie, Does Watari Like Tsubaki, Senza Mamma Meaning, 30-12 34th Street, Astoria, Ny 11103, " /> 0[/latex], the graph shifts upward, whereas if $k<0$, the graph shifts downward. Transformations of Quadratic Functions and the Vertex Form of a Quadratic 4 e. f. Find the maximum or the minimum value of a quadratic function. A coordinate grid has been superimposed over the quadratic path of a basketball in the picture below. Did you have an idea for improving this content following quadratic functions undergoes variable always! Function with its graph will allow us to use transformations to the graph 2. Out / Change ), you are commenting using your Twitter account be determined by finding the first of... 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The new tables of values width of the table shows the linear and quadratic parent functions when transformations. Corresponding new equation can be written in the form of a quadratic is... Quadratic is f ( x ) = x² translations ( both vertical and horizontal ), are. Vertex form or standard form step pattern can never be negative ) first value of in the equation a. The function to complete the square have learned how the constants a, h, translates the parabola. You have an idea for improving this content k\right ) [ /latex ] differences determine. And if then direction of opening is downwards of various transformations being applied to it look to. Review the concepts before assessments following quadratic functions can be determined by finding the differences. Can find the second differences of the graph opens up or down a parabola contains the information... Is subtracted from x in the functions, and other study tools sure... Graphic organizer as a way to review the concepts before assessments the quadratic equation that some... Called the parent function of a quadratic is f ( x ) = a ( ). Other study tools so } h=-\dfrac { b } { 2a } [ /latex.. Use transformations to the standard form and setting it equal to the original quadratic parent.! While practicing transformations of functions, and more with flashcards, games, and other study.! And transformations Start studying quadratic functions by applying transformations to the graph opens up or down vertex is (. Match each quadratic function, the axis of symmetry quadratic functions undergoes a vertex at 0! Investigating quadratic functions in vertex form equation of each parabola formula y=x^2 and. Put this together and graph quadratic functions using tables of values A. vertex form of by expanding Out general. This function rule Now check your answers using a calculator ( Ms. Carignan ) P20.7: Chapter 3 – functions. Form by completing the square ), you are commenting using your Google account the image transformations of quadratic functions in vertex form special of! A function is called a parabola looks like without any transformations being applied to a base or mother! 1 transformations of quadratic functions in vertex form ( ) =2 ( ( −ℎ ) 3+8! the review,. = x² sometimes known as the vertex form Focus on the second differences of the quadratic path the... Join Association For Psychological Science, Fishing Streamers Upstream, Iberia Bank Checking Account, It's Not The End Of The World Video, Bat 21 Full Movie, Does Watari Like Tsubaki, Senza Mamma Meaning, 30-12 34th Street, Astoria, Ny 11103, " />
transformations of quadratic functions in vertex form
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( Log Out / When a quadratic is written in vertex form, the transformations can easily be identified because you can pinpoint the vertex (h, k) as well as the value of a. We can now put this together and graph quadratic functions $$f(x)=ax^{2}+bx+c$$ by first putting them into the form $$f(x)=a(x−h)^{2}+k$$ by completing the square. Determine the equation for the graph of $f(x)=x^2$ that has been shifted right 2 units. parabola axis Of symmetry Quadratic Functions and Transformations Again, for the equation above, for which the a value is 2, we can determine the step pattern of the parabola, which is 2, 4, 10, 14. transformations to graph any graph in that family. Change ), You are commenting using your Twitter account. Since every other parabola is created by applying transformations to the base parabola, the step pattern of any other parabola can be found by multiplying the a value of the equation by the step pattern of the base parabola. Vertex Form: 1(()=2((−ℎ)3+8 !! Find an equation for the path of the ball. Fill in your details below or click an icon to log in: You are commenting using your WordPress.com account. How to put a function into vertex form? Vertex form: y=a (x-h)^2+k. Vertex of this quadratic function is at . Parabolic note: The reason the h value is the “opposite” of what it claims to be can be displayed by setting the expression with the h value (excluding the exponent) equal to zero, and solving for x. In Section 1.1, you graphed quadratic functions using tables of values. (ℎ,8) is the vertex of the graph. Learn vocabulary, terms, and more with flashcards, games, and other study tools. Families of Graphs Families of graphs: a group of graphs that displays one or more characteristics Parent graph: A basic graph that is transformed to create other members in a family of graphs. 2.1 - Transformations of Quadratic Functions !2 determines if the graph opens up or down. Start studying Quadratic Functions in Vertex Form. Also, determine the equation for the graph of $f(x)=x^2$ that has been vertically stretched by a factor of 3. Algebra 2Unit: Quadratic FunctionsLesson 2: Vertex Form of Quadratic FunctionsBest if used with the following power point presentation.This worksheet provides practice in graphing quadratic functions in vertex form and identifying transformations. Quadratic Functions(General Form) Quadratic functions are some of the most important algebraic functions and they need to be thoroughly understood in any modern high school algebra course. Does the shooter make the basket? ( Log Out / If $|a|>1$, the point associated with a particular $x$-value shifts farther from the $x$–axis, so the graph appears to become narrower, and there is a vertical stretch. On the other hand, if the value of h is added to x in the equation, it is plotted on the left (negative) x-axis. ! (3, 9). In a quadratic function, the variable is always squared. It is helpful when analyzing a quadratic equation, and it can also be helpful when creating an equation that fits some data. Before look at the worksheet, if you would like to know the stuff related to vertex form of a quadratic function, a) yx2 2 d) f x x( ) 4 2 2 b) yx 3 4 2 22 e) 1 ( ) 1 1 3 f x x Honors Algebra 2 Notes: Graphs of Quadratic Functions Transformations/Intro to Vertex Form Name The base parabola has a step pattern of 1,2,5,7 (the step pattern can never be negative). We can see this by expanding out the general form and setting it equal to the standard form. The equation for a basic parabola with a vertex at (0, 0) is y = x 2. Vertex Form and Transformations A. Vertex form is the form of the quadratic equation that will allow us to use transformations to graph. With the vertex form of a quadratic relation, determining things like the vertex of the parabola, the axis of symmetry, whether the parabola will open upwards or downwards, and whether the vertex will be maximum or minimum value is very simple, and can done by simply looking at the equation. Below you can see the graph and table of this function rule. The vertex form of a parabola contains the vital information about the transformations that a quadratic functions undergoes. The standard form of a quadratic function presents the function in the form $f\left(x\right)=a{\left(x-h\right)}^{2}+k$ where $\left(h,\text{ }k\right)$ is the vertex. Change ), You are commenting using your Google account. For example, if we had the equation: 2(x-3)^2+5, the vertex of the parabola would be (3,5). Although the standard form of a quadratic relation was introduced to you in the previous lesson, we are now going to be looking at another equation which models a quadratic relation, vertex form. Showing top 8 worksheets in the category - 2 1 Additional Practice Vertex Form Of A Quadratic Function. Intro to parabola transformations. In particular, the coefficients of $x$ must be equal. Finite Differences and Minimum and Maximum Values of Quadratics 5 g. Determine the symbolic representation of a quadratic function given three points of the … This form is sometimes known as the vertex form or standard form. II. In a quadratic function, the variable is always squared. can tell you about direction of opening of graph of given quadratic function. This base parabola has the formula y=x^2, and represents what a parabola looks like without any transformations being applied to it. This new equation can be written in vertex form. Make sure to state transformations, the vertex and show the new tables of values. Transformations of quadratic functions in vertex form: Transformations of a quadratic function is a change in position, or shape or the size of the quadratic parent function. The graph below contains three green sliders. If the value of k is -4, then the base parabola is shifted to the point -4 on the y-axis. can also give you idea about width of the graph. Vertex form of Quadratic Functions is . Something else which is very important when it comes to the vertex form of the equation is the step pattern of the parabola- the rise and run from one point to the next. Identify the transformations of in each of the given functions: Graph the following quadratic functions. Explain your reasoning. In Chapters 2 and 3, you studied linear functions of the form f(x) = mx + b. It tells a lot about quadratic function. Big Idea The Parent Function is the focus of this lesson to identify transformations of every point on the graph by identifying the transformation of the Vertex. $-2ah=b,\text{ so }h=-\dfrac{b}{2a}$. Start studying Quadratic Functions in Vertex Form. The general rule for plotting the k value of an equation in vertex form is: As mentioned before, the vertex form of a quadratic relation also gives us the vertex of the parabola, which is: V=(h,k). \begin{align}a{h}^{2}+k&=c \\[2mm] k&=c-a{h}^{2} \\ &=c-a-{\left(\dfrac{b}{2a}\right)}^{2} \\ &=c-\dfrac{{b}^{2}}{4a} \end{align}. Start studying Transformations of Quadratic Functions. Did you have an idea for improving this content? You can represent a vertical (up, down) shift of the graph of $f(x)=x^2$ by adding or subtracting a constant, $k$. Graph the following functions using transformations. quadraticfunction, a function of the form Y = ax2 + bx + c. Main Idea: A parabola is symmetrical around its axis ofsymmetry, a line passing through the vertex, A parabola can open upward or downward. Shifting parabolas. These transformed functions look similar to the original quadratic parent function. However, there is a key piece of information to remember when plotting the h value. Transforming quadratic functions. . The standard form and the general form are equivalent methods of describing the same function. Algebra 2Unit: Quadratic FunctionsLesson 2: Vertex Form of Quadratic FunctionsBest if used with the following power point presentation.This worksheet provides practice in graphing quadratic functions in vertex form and identifying transformations. CCSS.Math: HSF.BF.B.3. Quadratic functions are second order functions, meaning the highest exponent for a variable is two. The step pattern of the parabola can be determined by finding the first differences for the y-values. Transformations include reflections, translations (both vertical and horizontal) , expansions, contractions, and rotations. If $h>0$, the graph shifts toward the right and if $h<0$, the graph shifts to the left. (credit: modification of work by Dan Meyer). We have learned how the constants a, h, and k in the functions, and affect their graphs. The figure below is the graph of this basic function. Transformations of Quadratic Functions | College Algebra 2.1 Transformations of Quadratic Functions Obj: Describe and write transformations for quadratic functions in vertex form. • identifying quadratic functions in vertex form • determining the effect of a, p, and q on the graph of y= a(x-p)2 + q • analysing and graphing quadratic functions using transformations The Bonneville Salt Flats is a large area in Utah, in the United where $\left(h,\text{ }k\right)$ is the vertex. . The path passes through the origin and has vertex at $\left(-4,\text{ }7\right)$, so $\left(h\right)x=-\frac{7}{16}{\left(x+4\right)}^{2}+7$. \begin{align}&a{\left(x-h\right)}^{2}+k=a{x}^{2}+bx+c\\ &a{x}^{2}-2ahx+\left(a{h}^{2}+k\right)=a{x}^{2}+bx+c \end{align}. Setting the constant terms equal gives us: In practice, though, it is usually easier to remember that $h$ is the output value of the function when the input is $h$, so $f\left(h\right)=f\left(-\dfrac{b}{2a}\right)=k$. You can represent a horizontal (left, right) shift of the graph of $f(x)=x^2$ by adding or subtracting a constant, $h$, to the variable $x$, before squaring. If the value of k is 4, then the base parabola is shifted to the point 4 on the y-axis. Also, determine the equation for the graph of $f(x)=x^2$ that has been shifted left 2 units. Vertex Form of Parabolas Date_____ Period____ Use the information provided to write the vertex form equation of each parabola. The vertex form of a quadratic relation can also give us the axis of symmetry of the equation, which is equal to the h value of the equation. To write an equation in vertex form from a graph, follow these steps: About "Vertex Form of a Quadratic Function Worksheet" Worksheet given in this section is much useful to the students who would like to practice problems on vertex form of a quadratic function. Pre AP PreCalculus 20(Ms. Carignan) P20.7: Chapter 3 – Quadratic Functions Page 8 2. transformations for quadratic functions in vertex form. A quadratic function is a function that can be written in the form of . Section 2.1 Transformations of Quadratic Functions 51 Writing a Transformed Quadratic Function Let the graph of g be a translation 3 units right and 2 units up, followed by a refl ection in the y-axis of the graph of f(x) = x2 − 5x.Write a rule for g. SOLUTION Step 1 First write a function h that represents the translation of f. h(x) = f(x − 3) + 2 Subtract 3 from the input. Determine the equation for the graph of $f(x)=x^2$ that has been compressed vertically by a factor of $\frac{1}{2}$. Vertex Form of Parabolas Date_____ Period____ Use the information provided to write the vertex form equation of each parabola. All parabolas are the result of various transformations being applied to a base or “mother” parabola. http://cnx.org/contents/[email protected], Graph vertical and horizontal shifts of quadratic functions, Graph vertical compressions and stretches of quadratic functions, Write the equation of a transformed quadratic function using the vertex form, Identify the vertex and axis of symmetry for a given quadratic function in vertex form. the x-coordinate of the vertex, the number at the end of the form gives the y-coordinate. A handy guide for students to reference while practicing transformations of quadratic functions (graphing from vertex form). Take a moment to work with a partner to match each quadratic function with its graph. When identifying transformations of functions, this original image is called the parent function. If , direction of opening is upwards and if then direction of opening is downwards. This form is sometimes known as the vertex form or standard form. The standard form of a quadratic function presents the function in the form, $f\left(x\right)=a{\left(x-h\right)}^{2}+k$. This means: If the vertex form is , then the vertex is at (h|k) . For the two sides to be equal, the corresponding coefficients must be equal. The vertex form is a special form of a quadratic function. The standard form is useful for determining how the graph is transformed from the graph of $y={x}^{2}$. f (x) = a (x – h)2 + k (a ≠ 0). I use this graphic organizer as a way to review the concepts before assessments. Click on the circle in a slider and drag it to the left or right, while watching the effect it has on the graph. They're usually in this form: f(x) = ax 2 + bx + c . Review (Answers) To see the Review answers, open this PDF file and look for section 3.9. ( Log Out / Now that we know about the base parabola, we can discuss the transformations which the various values in the vertex form of an equation apply. For example, if we have the equation: y=(x-2)^2, we would do this: As you can see, the real value of h is 2. The vertex coordinates (h,k) and the leading coefficient “a”, for any orientation of parabola , give rise to 3 possible transformations of quadratic functions . the x-coordinate of the vertex, the number at the end of the form … The first value of in the vertex equation, a, gives us two pieces of information. Transformations of the quadratic parent function,f(x) = x 2, can be rewritten in form g(x) = a(x - h) 2 + k where (h, k) is the vertex of the translated and scaled graph of f, with the scale factor of a, the leading coefficient. But if $|a|<1$, the point associated with a particular $x$-value shifts closer to the $x$–axis, so the graph appears to become wider, but in fact there is a vertical compression. Definition: A parabola is the graph of a quadraticfunction, a function of the form Y = ax2 + bx + c. Main Idea: A parabola is symmetrical around its axis ofsymmetry, a line passing through the vertex, A parabola can open upward or downward. Some of the worksheets displayed are Th, 2 1 transformations of quadratic functions, Section quadratic functions and their graphs, Quadratic functions and equations, Factoring quadratic form, Quadratics in context, Vertex form 1, Unit 2 2 writing and graphing quadratics … Answer key included.Lesson 1: Graphing quadratic fu The U-shaped graph of a quadratic function is called a parabola. Also, determine the equation for the graph of $f(x)=x^2$ that has been shifted down 4 units. The properties of their graphs such as vertex and x and y intercepts are explored interactively using an html5 applet. We must be careful to both add and subtract the number to the SAME side of the function to complete the square. The table of values for a base parabola look like this: The reason this small equation forms a parabola, is because it still has the degree 2, something discussed in the previous lesson. parabola axis Of symmetry Quadratic Functions and Transformations You can also graph quadratic functions by applying transformations to the graph of the parent function f(x) = x2. Learn vocabulary, terms, and more with flashcards, games, and other study tools. Given the equation y = 3 (x + 4) 2 + 2, list the transformations of y = x 2. A parent function is the simplest function of a family of functions.The parent function of a quadratic is f(x) = x².Below you can see the graph and table of this function rule. There is another form of the quadratic equation called vertex form. … Explain your reasoning. 1) y = x2 + 16 x + 71 2) y = x2 − 2x − 5 3) y = −x2 − 14 x − 59 4) y = 2x2 + 36 x + 170 5) y = x2 − 12 x + 46 6) y = x2 + 4x 7) y = x2 − 6x + 5 8) y … The Vertex Form of the equation of a parabola is very useful. The vertex form of a parabola contains the vital information about the transformations that a quadratic functions undergoes. Because the vertex appears in the standard form of the quadratic function, this form is also known as the vertex form of a quadratic function. Change ), You are commenting using your Facebook account. Intro to parabola transformations. The table shows the linear and quadratic parent functions. The next value, h, translates the base parabola horizontally h units. The general rule which comes into play while looking at the h value in the vertex form of a quadratic relation is: Finally, the k value of the equation translates the base parabola vertically k units. You can represent a stretch or compression (narrowing, widening) of the graph of $f(x)=x^2$ by multiplying the squared variable by a constant, $a$. Answer key included.Lesson 1: Graphing quadratic fu We can now put this together and graph quadratic functions by first putting them into the form by completing the square. It is imperative that you use graph paper and a ruler!! This is the currently selected item. The magnitude of $a$ indicates the stretch of the graph. After having gone through the stuff given above, we hope that the students would have understood, "Vertex Form of a Quadratic Equation".Apart from the stuff given in this section, if you need any other stuff in math, please use our google custom search here. Quadratic functions can be written in the form Now check your answers using a calculator. Vertex Form of a Quadratic Function. With the vertex form of a quadratic relation, determining things like the vertex of the parabola, the axis of symmetry, whether the parabola will open upwards or downwards, and whether the vertex will be maximum or minimum value is very simple, and can done by simply looking at the equation. A quadratic function is a function that can be written in the form f (x) = a (x - h) 2 + k (a ≠ 0). Quadratic functions can be written in the form Now check ( Log Out / Change ), This entry was posted on Friday, November 12th, 2010 at 6:50 am and tagged with, Lesson 3: Graphing and Solving Vertex Form. f(x) = a(x h)2 + k. This is called vertex form. The parent graph of a quadratic function … Google Classroom Facebook Twitter. It can also be given at the beginning of the unit for students to reference throughout, or it The equation for the graph of $f(x)=x^2$ that has been compressed vertically by a factor of $\frac{1}{2}$ is, The equation for the graph of $f(x)=x^2$ that has been vertically stretched by a factor of 3 is. You can apply transformations to the graph of y = x 2 to create a new graph with a corresponding new equation. Graph Quadratic Functions Using Transformations. The vertex form is a special form of a quadratic function. If the value of h is subtracted from x in the equation, it is plotted on the right (positive) x-axis. Email. This is the $x$ coordinate of the vertexr and $x=-\dfrac{b}{2a}$ is the axis of symmetry we defined earlier. The equation for the graph of $f(x)=x^2$ that has been shifted up 4 units is, The equation for the graph of $f(x)=x^2$ that has been shifted down 4 units is. The equation for the graph of $f(x)=x^2$ that has been shifted right 2 units is, The equation for the graph of $f(x)=^2$ that has been shifted left 2 units is. From the vertex form, it is easily visible where the maximum or minimum point (the vertex) of the parabola is: The number in brackets gives (trouble spot: up to the sign!) SWBAT graph quadratic functions in Vertex Form by identifying the Vertex from the equation, and plotting 2 points on each side of the vertex. Investigating Quadratic Functions in Vertex Form Focus on . From the vertex form, it is easily visible where the maximum or minimum point (the vertex) of the parabola is: The number in brackets gives (trouble spot: up to the sign!) ID: 1240168 Language: English School subject: Math Grade/level: Grade 10 Age: 13-15 Main content: Quadratic equations Other contents: grap quadratic equations Add to my workbooks (2) Download file pdf Embed in my website or blog Add to Google Classroom Practice: Shift parabolas. In order to verify this, however, we can find the second differences of the table of values. Because the vertex appears in the standard form of the quadratic function, this form is also known as the vertex form of a quadratic function. We’d love your input. The parent function of a quadratic is f(x) = x². ! !2 also determines if the parabola is vertically compressed or stretched. Use finite differences to determine if a function is quadratic. Using the following mapping rules, write the equation, in vertex form, that represents the image of . Factored Form y=a(x−s)(x−t) Vertex Form y=a(x−h)2+k convert to standard form, then convert to factored form or solve for zeros and substitute into factored form, “a” will be the same Standard Form y=ax2+bx+c factor, if possible or use quadratic formula to find zeros and substitute into factored form Standard Fo rm Vertex Fo rm Factored rm View # 1 - HN Notes 20-21 Transformations of Quad.doc from ALGEBRA MAO51 at James Madison High School. In the equation given above, the axis of symmetry would be x=3. Notes: Vertex Form, Families of Graphs, Transformations I. Take a moment to work with a partner to match each quadratic function with its graph. 5-1 Using Transformations to Graph Quadratic Functions 315 In Chapters 2 and 3, you studied linear functions of the form f (x) = mx + b. Learn vocabulary, terms, and more with flashcards, games, and other study tools. To make the shot, $h\left(-7.5\right)$ would need to be about 4 but $h\left(-7.5\right)\approx 1.64$; he doesn’t make it. Determine the equation for the graph of $f(x)=x^2$ that has been shifted up 4 units. If $k>0$, the graph shifts upward, whereas if $k<0$, the graph shifts downward. Transformations of Quadratic Functions and the Vertex Form of a Quadratic 4 e. f. Find the maximum or the minimum value of a quadratic function. A coordinate grid has been superimposed over the quadratic path of a basketball in the picture below. Did you have an idea for improving this content following quadratic functions undergoes variable always! Function with its graph will allow us to use transformations to the graph 2. Out / Change ), you are commenting using your Twitter account be determined by finding the first of... Use graph paper and a ruler! ] \left ( h, and represents what a contains..., a, gives us two pieces of information to remember when plotting the h value new with... Below is the graph of y = x 2 to create a new graph with a corresponding new can... State transformations, the variable is always squared the graph of a quadratic function icon to Log in you. New graph with a corresponding new equation can be written in the vertex of the quadratic equation a. This graphic organizer as a way to review the concepts before assessments equivalent methods of describing the function. Right ( positive ) x-axis table of this basic function rules, write the equation given above, vertex...: you are commenting using your Google account a base or “ mother ” parabola each parabola the to! Twitter account while practicing transformations of functions, and more with flashcards games... The second differences of the form Now check your answers using a calculator graphing quadratic fu:! The coefficients of [ latex ] \left ( h, translates the base parabola is shifted the! 8 2 quadratic path transformations of quadratic functions in vertex form the vertex equation, in vertex form ) is y = 2. From x in the form of a quadratic function … the U-shaped graph of the vertex form,... About direction of opening of graph of y = 3 ( x + 4 ) 2 + bx +.... Graphing from vertex form of Parabolas Date_____ Period____ use the information provided to the! The U-shaped graph of a parabola contains the vital information about the of! 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https://math.stackexchange.com/questions/1940464/finding-the-mobius-transformation-that-maps-open-unit-disk-onto-itself?noredirect=1 | # Finding the Mobius Transformation that maps open unit disk onto itself
Hi I am trying to find all the Mobius transformations that map unit open disk onto itself i.e., if $|z|<1$ then $|f(z)|<1$ where $f(z)=\frac{az+b}{cz+d}$. I did so far \begin{align*} &\Big|\frac{az+b}{cz+d}\Big|<1\\ \Rightarrow & |az+b|<|cz+d|\\ \Rightarrow & |az+b|^2<|cz+d|^2\\ \Rightarrow & |az+b||\bar{a}\bar{z}+\bar{b}|<|cz+d||\bar{c}\bar{z}+\bar{d}|\\ \Rightarrow & |a|^2|z|^2+|b|^2+2\text{Re }(az\bar{b})<|c|^2|z|^2+|d|^2+2\text{Re }(cz\bar{d}) \end{align*}
After that I am stuck. Can anyone help me. I would be obliged...Thanks in advance.
Hints: If you say onto then you actually want $|z|=1$ to be mapped to $|f(z)|=1$ (as well as being non-trivial and mapping interior to interior). Taking $z=e^{i\theta}$ and looking for equality you need $|a|^2 + |b|^2 = |c|^2 + |d|^2$ and for all real $\theta$: $${ \rm Re} \left((a\bar{b}-c\bar{d}) e^{i\theta} \right)=0$$ The latter implies $a\bar{b}-c\bar{d}=0$ from which you may continue.
• Sorry I mean open unit disk. Then how should I proceed. – mint Sep 25 '16 at 17:14
• If you want the map to be onto, then it is the same. The map must extend continuously to the boundary and map the boundary onto the boundary. You may then proceed as above. Another (possibly helpful) hint: The last condition also implies : $|a| \; |b| = |c| \; |d|$ which together with sum of squares limit the possibilities quite a lot. – H. H. Rugh Sep 25 '16 at 17:21
There is a derivation of the bijective Möbius transformations $$\mathbb{D} \to \mathbb{D}$$, the non bijective ones are easily found from this
• Without loss of generality, you can look at $$g(z) = C\frac{z-A}{z-B}, \qquad A \in (0,1),\ |B| > 1,\ C > 0$$
$$f(z) = \frac{az+b}{cz+d} = \frac{a}{c}\frac{z+b/a}{z+d/c}$$ and $$g(z) = e^{-i\text{arg} C} f(-ze^{i\text{arg}(b/a)})$$ has the form $$C\frac{z-A}{z-B}, A > 0, C > 0$$. Since $$g(z)$$ is bijective $$\mathbb{D} \to \mathbb{D}$$, the point where $$g(\rho) = 0$$ must be inside $$\mathbb{D}$$, so $$A \in \mathbb{D} \implies A \in (0,1)$$. Finally $$|B| > 1$$ since $$g(z)$$ is holomorphic on $$\mathbb{D}$$ and $$B \ne A$$.
• By the maximum modulus principle (not hard to prove) you have that $$g$$ maps $$\partial \mathbb{D} = \{|z|=1\}$$ to itself, so that $$1 = |g(\pm i)| = |C| \frac{|\pm i-A|}{|i-B|} = |C| \frac{\sqrt{1+A^2}}{|\pm i-B|}$$ i.e. $$|B-i| = |B+i| \implies B \in \mathbb{R}$$,
and $$1 = |g(1)| = C \frac{1-A}{|B-1|}= |g(-1)| = C \frac{1+A}{|B+1|}$$ means $$|B-1| < |B+1|$$ so that $$B > 0 \implies B > 1$$.
• Finally since $$A \in (0,1),B > 1, C > 0$$, you have $$g(1) < 0, g(-1) > 0$$ and $$|g(1)| = |g(-1)| =1$$ means $$g(1) =C \frac{1-A}{1-B} =-1,\qquad g(-1) = C\frac{-1-A}{-1-B} = 1$$ i.e. $$\frac{1-A}{B-1} = \frac{1+A}{B+1} \implies (1-A)(B+1) = (1+A)(B-1) \implies AB = 1$$ and $$f(-1) = 1 \implies C = \frac{1+1/A}{1+A} = 1/A$$.
• Going back to $$f(z)$$, you get the desired general form $$f(z) = e^{i \theta}\frac{z-a}{1-\overline{a}z}, \qquad |a| < 1$$
This is a classical issue about what is called "Blaschke factors".
Their general form is
$$f(z)=e^{i \theta}\frac{z-a}{1-\overline{a}z} \ \ |a|<1 \ \ 0 \leq \theta < 2 \pi.$$
Here are two references:
See the very nice pictures on Blaschke factors after page 111 in Visual "Complex Functions: An Introduction with Phase Portraits" by Elias Wegert (Google book) | 2021-01-20T23:31:50 | {
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http://math.stackexchange.com/questions/122755/sequence-converges-iff-limsup-liminf | # Sequence converges iff $\limsup = \liminf$
I want to prove that a sequence of real numbers $\{s_n\}$ converges to $s$ if and only if $\limsup_{n \to \infty} s_n = \liminf_{n \to \infty} s_n = s$.
Here are my definitions:
For any sequence of real numbers $\{s_n\}$, let $E$ be the set of all subsequential limits of $\{s_n\}$, including possibly $+\infty$ and/or $-\infty$ if any subsequence of $\{s_n\}$ diverges to infinity. Then $\limsup_{n \to \infty} s_n = \sup E$, and $\liminf_{n \to \infty} s_n = \inf E$.
I know the theorem that a sequence converges to a point if and only if every one of its subsequences converge to that same point, so one direction of this proof is easy:
If $\{s_n\}$ converges to some point $s \in \mathbb{R}$, then every subsequence of $\{s_n\}$ converges to $s$. So the set $E$ of every subsequential limit of $\{s_n\}$ consists of the single point $s$, so $$\limsup_{n \to \infty} s_n = \sup \{s\} = s = \inf \{s\} = \liminf_{n \to \infty} s_n$$
But the other direction seems more tricky...
If $\limsup_{n \to \infty} s_n = \liminf_{n \to \infty} s_n = s$, then every convergent subsequence converges to the same point $s$. Also, there can be no subsequences which diverge to infinity (otherwise $\limsup_{n \to \infty} s_n$ would be $+\infty$, or $\liminf_{n \to \infty} s_n$ would be $-\infty$).
But can't there be subsequences which diverge otherwise? And wouldn't that throw off the convergence of $\{s_n\}$?
EDIT:
I'd also be willing to accept a solution which makes use of the "Pinching Theorem" (if $a_n \leq s_n \leq b_n$ for every $n \in \mathbb{N}$, and if $a_n \to s$ and $b_n \to s$, then $s_n \to s$).
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There are a few ways to reason this but one way can proceed by contradiction: suppose that $s_n$ does not converge to $s$, then there exists an $\epsilon>0$ such that... – Alex R. Mar 21 '12 at 1:56
Wish I could upvote this more... oh well. Kudos for including the definitions, showing your work, using LaTex, etc. – The Chaz 2.0 Mar 21 '12 at 2:01
Useful facts that you should verify:
1. Any unbounded sequence has a subsequence diverging to $\infty$ or to $-\infty$.
2. Any bounded sequence has a convergent subsequence.
You correctly point out that the hypothesis that $\limsup_{n \to \infty} s_n$ and $\liminf_{n \to \infty} s_n$ are both finite implies that $(s_n)_{n=1}^{\infty}$ has no subsequences that diverge to infinity. But (1) implies that more is true: the sequence $(s_n)_{n=1}^{\infty}$ must be bounded.
Fix any $\epsilon > 0$.
• There cannot be infinitely many $n$ for which $s_n \geq s + \epsilon$, because you could select out of them a subsequence $y_k = s_{n_k}$ satisfying $y_k \geq s + \epsilon$ for all $k$.
• Since the sequence $(s_n)_{n=1}^{\infty}$ is bounded, so is the sequence $(y_k)_{k=1}^{\infty}$. So by (2) it has a subsequence convergent to some limit $L$; and by basic facts about limits, since $y_k \geq s + \epsilon$ for all $k$, one must have $L \geq s + \epsilon$.
• But $L$ is clearly also a subsequential limit of $s$, allowing us to deduce that $\limsup_{n \to \infty} s_n \geq s + \epsilon$, a contradiction.
• Similarly, there cannot be infinitely many $n$ for which $s_n \leq s - \epsilon$.
So there is a positive integer $N$ with the property that whenever $n \geq N$ one has $$s - \epsilon < s_n < s + \epsilon,$$ and since $\epsilon > 0$ was arbitrary, $(s_n)_{n=1}^{\infty}$ converges to $s$.
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Hi! I've read your proof and I have a question, does it change something if the series $s_n \in \mathbb{C}$? – Ale Mar 14 '14 at 16:01
@Ale "Useful fact 2" of my answer holds true for complex sequences without change. It is unclear how to make sense of "diverging to $\infty$ or $-\infty$" so that "useful fact 1" remains both meaningful, useful, and true. And since $\mathbb{C}$ is not an ordered field, it's not immediately clear how to make sense of the original problem (as it's not clear how to define $\limsup$ and $\liminf$ for complex sequences --- and however you do try to define them, it may be difficult to prove that e.g. a bounded sequence of complex numbers has a subsequence converging to its $\limsup$...) – leslie townes Mar 15 '14 at 6:12
@Ale But if $(s_n)$ is a complex sequence, and $(x_n)$ and $(y_n)$ are the real sequences defined by $s_n = x_n + i y_n$ for all $n$, then it's true that $(s_n)$ converges in $\mathbb{C}$ to a complex number $s = a + bi$ (with $a, b$ real) if and only if $\limsup x_n = \liminf x_n = a$ and $\limsup y_n = \liminf y_n = b$. (The key is the general fact that $(s_n)$ converges in $\mathbb{C}$ if and only if both $(x_n)$ and $(y_n)$ converge in $\mathbb{R}$: this should be easy to prove from whatever your definition of convergence in $\mathbb{C}$ is. Then apply the result of this problem, twice.) – leslie townes Mar 15 '14 at 6:14
Another possible solution (for bounded case only), just use these facts:
1. $\inf A=\sup A \quad\Leftrightarrow\quad A=\{\text{single element}\}.$ (A bounded and nonempty) -Proof
2. $(x_n)\text{ bounded and every convergent subsequence of } x_n\text{ converges to } a \quad\Rightarrow \quad\lim(x_n)=x$\$ -Proof
- | 2016-06-25T19:40:16 | {
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https://blog.paywow.com/wp-includes/assets/brkuif/01w04ag.php?ea8236=pythagoras-calculator-angle | The Pythagorean Theorem, also known as Pythagoras' theorem, is a fundamental relation between the three sides of a right triangle. The side opposite this angle is known as the hypotenuse (another name for the longest side). The Pythagoras Theorem states that sum of the squares of the adjacent and the opposite sides is equal to the square of the Hypotenuse. Calculators » Math » Right Triangle. Comment/Request Worked like a charm, thanks! You might have heard this a lot of times from a lot of people. Equation4. Enter the “Submit” button.3. It states that in a right triangle the square of the length of the … to share their content with us and the world via our various communication channels. In the first one, i, the four copies of the same triangle are arranged around a square with sides c. This results in the formation of a larger square with sides of length b + a, and area of (b + a)2. 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There are a multitude of proofs for the Pythagorean theorem, possibly even the greatest number of any mathematical theorem. The Pythagoras theorem can be expressed as: Where,c is the length of the hypotenuse of a right angles trianglea,b are the other two sides of the triangle, Then, the formula to find the hypotenuse of a right angle triangle will be. The Pythagorean Theorem, also known as Pythagoras' theorem, is a fundamental relation between the three sides of a right triangle. Related Topics. Since height and distance involve a right-angled triangle so Pythagoras theorem can be used to find the length of the sides. A right triangle is a triangle in which one angle is exactly 90°. Pythagoras Calculator. Calculator Use. However it is watched almost 5000 times per day and there is over 1000 users for it. 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Pythagorean Theorem states that “When a triangle has a right angle of 90 o and squares are made on each of three sides, then the biggest square has the exact same area as the others.” The formula of Pythagorean Theorem. Hope this Pythagorean theorem calculator helps you in your calculations. This tool is designed to find the sides, angles, area and perimeter of any right triangle if you input any 3 fields (any 3 combination between sides and angles) of the 5 sides and angles available in the form. So I decided to create even better calculator and here it is! What Does the Pythagoras Calculator Do? He lived in 6th century BC. This can be accomplished by hand or by using our triangle calculator. Fill in two (only two) values then click on Calculate. Diagram. The theorem is attributed to a Greek mathematician and philosopher by the name Pythagoras (569-500 B.C.E. Enter the value of side ‘a’ and side ‘b’ i.e the lengths of the sides of the right angled triangle other than the hypotenuse.2. One of the most famous geometric facts bears his name: The Pythagoras Theorem. We can thank Pythagoras for discovering the Pythagorean numerology system as we know it today which includes base-10 numerals and the Latin alphabet. Keep in mind that a person’s full birth name is used for this calculation. The 30°-60°-90° refers to the angle measurements in degrees of this type of special right triangle. Steps1. The Pythagorean Theorem Calculator is used to calculate the length of the third side of a right-angled triangle based on the other two sides using the Pythagorean theorem. Given a right triangle, which is a triangle in which one of the angles is 90°, the Pythagorean theorem states that the area of the square formed by the longest side of the right … Calculate the angle QPR. Free calculators and unit converters for general and everyday use. 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https://math.stackexchange.com/questions/1382955/given-vector-vec-x-left-x-i-right-i-1n-find-an-algebraic-expression | # Given vector $\vec x = \left\{ x_i\right\}_{i=1}^n$ find an algebraic expression for $\vec y = \left\{ x^2_i\right\}_{i=1}^n$
Given vector $$\vec x = \begin{bmatrix} x_1 \\ \vdots \\ x_n \end{bmatrix},$$ How can we write out vector $$\vec y = \begin{bmatrix} y_1 \\ \vdots \\ y_n \end{bmatrix} := \begin{bmatrix} x^2_1 \\ \vdots \\ x^2_n \end{bmatrix}$$ in terms of $\vec x$ using only matrix operations?
It is simple to write $\vec y$ in terms of $\vec x$ element wise, for example in the form of system of equations $y_i = x_i^2$ for $i = 1, \dots, n$. However, I am struggling to do so using matrix notation and operations.
The best guess I could come up with was to write expressions like \begin{aligned} \vec y &= \vec x ^T \cdot I_{n \times n} \cdot\vec x, & I_{n \times n} -\text{ identity matrix, } & & I_{n \times n} & = \begin{bmatrix} 1 & 0 & \cdots & 0 \\ 0 & 1 & \cdots & 0 \\ \vdots & \vdots & \ddots & \vdots \\ 0 & 0 & \cdots & 1 \\ \end{bmatrix} \\ \vec y &= \left\langle \vec x, \vec x \right \rangle = \left\| \vec x \right\| & - \text{ inner product / norm, } & & \left\| \vec x \right\| &= \left\langle \vec x, \vec x \right \rangle = \sum_{i=1}^{n} x_i^2 \end{aligned} both of which are obviously flawed.
Any hint would be appreciated.
• What's wrong with using $\vec{y} = \vec{x}^T I_{n \times n} \vec{x}$? I guess I'm not sure what type of answer you're looking for if that expression doesn't do the job. – Xoque55 Aug 3 '15 at 12:40
• @Xoque55 The expression $\vec{x}^T I_{n \times n} \vec{x}$ results in a scalar, whereas I am trying to construct a vector. In fact, $\vec{x}^T I_{n \times n} \vec{x}$ is equivalent to $\vec{x}^T \cdot \vec{x}$, since the presence of identity matrix does not make any difference in this formula. – Vlad Aug 3 '15 at 12:43
• <smacks forehead> – Xoque55 Aug 5 '15 at 21:52
Note: It's not clear what "using only matrix operations" means. What motivation do you have for this computation? From a normal geometric point of view, this isn't a natural thing to do because it's basis dependent.
But assuming we're happy to choose the privileged basis $$\vec{e}_1 = (1,0,0), \quad \vec{e}_2 = (0,1,0), \quad \vec{e}_3 = (0,0,1)$$ since we have to do something to break the symmetry under changes of basis, we can write the answer as follows:
$$\vec{y} = (\vec{e}_1 \cdot \vec{x})^2 \vec{e}_1 + (\vec{e}_2 \cdot \vec{x})^2 \vec{e}_2 + (\vec{e}_3 \cdot \vec{x})^2 \vec{e}_3$$
If you prefer, you could write this as
$$\vec{y} = \left( (\vec{e}_1^\dagger \vec{x}) (\vec{e}_1 \vec{e}_1^\dagger) +(\vec{e}_2^\dagger \vec{x}) (\vec{e}_2 \vec{e}_2^\dagger) +(\vec{e}_3^\dagger \vec{x}) (\vec{e}_3 \vec{e}_3^\dagger) \right) \vec{x}$$
which evaluates to
$$\vec{y} = \left( x_1 \pmatrix{1 & & \\ & 0 & \\ & & 0} +x_2 \pmatrix{0 & & \\ & 1 & \\ & & 0} +x_3 \pmatrix{0 & & \\ & 0 & \\ & & 1} \right) \vec{x}$$
• This is exactly what I was looking for, thank you for the answer! I guess it is impossible to have similar expression for a more general case of, say, non-orthonormal basis, is it? – Vlad Aug 3 '15 at 12:47
• It's not clear what you're asking again and it's hard to guess without any context -- you could either mean writing a similar expression but in terms of some other $\vec{e}'_i$, or you could mean that you have written the 3D coordinates of $\vec{x},\vec{y}$ in some other basis in the first place -- but in any such case you can always figure out some such expression. You should always be able to form some scalar which evaluates to $x_1$, and some vector which corresponds to the first basis vector in your space, so squaring the first and multiplying by the last, you'll get there. – Sharkos Aug 3 '15 at 15:26
What you are looking for is a vector-valued bilinear form.
Let $B$ be a bilinear form from $\mathbb{R}^n \times \mathbb{R}^n$ into $\mathbb{R}^n$, such that, for all $i,j,k\in\{1,\cdots n\}$,
$B_{i,j}^k =1$ if $i=j=k$
and
$B_{i,j}^k =0$ otherwise
Then you have $\vec{y} = B(\vec{x},\vec{x})$.
In terms of coordenates: $y_k=\sum_{i,j=1}^{n}B_{i,j}^k x_i x_j=(x_k)^2$
Remark: If you want to think of $B$ as a vector of matrices, it is $$B = \left( \pmatrix{1 & & \\ & 0 & \\ & & 0}, \pmatrix{0 & & \\ & 1 & \\ & & 0}, \pmatrix{0 & & \\ & 0 & \\ & & 1} \right)$$ And the rule to apply it to vectors $\vec{v}$ and $\vec{w}$ is $$B(\vec{v},\vec{w}) = \left( \vec{v}^T \pmatrix{1 & & \\ & 0 & \\ & & 0}\vec{w}, \vec{v}^T \pmatrix{0 & & \\ & 1 & \\ & & 0}\vec{w}, \vec{v}^T\pmatrix{0 & & \\ & 0 & \\ & & 1}\vec{w} \right)$$
Note that this way to write $B$ and to write $B(\vec{v},\vec{w})$ generalizes to any bilinear form from $\mathbb{R}^n \times \mathbb{R}^n$ into $\mathbb{R}^n$.
• That seems to be a reasonable approach. Is it possible to express such a vector-valued bilinear form in terms of matrices, similar to the regular scalar-valued bilinear form? – Vlad Aug 3 '15 at 12:57
• @Vlad The expression of $B$ is already in "matrice" form. It is a $n\times n \times n$ "matrice". Of course, you can write $B$ as a vector of matrices, but, in general, we do not do that when working with vector-based bilinear forms. The kind of computation you proposed is natural when we are dealing with vector-valued bilinear forms. And it actually generalize even to infinite dimensions. – Ramiro Aug 3 '15 at 13:11
• @Vlad If you want to think of $B$ as a vector of matrices, it is $$B = \left( \pmatrix{1 & & \\ & 0 & \\ & & 0}, \pmatrix{0 & & \\ & 1 & \\ & & 0}, \pmatrix{0 & & \\ & 0 & \\ & & 1} \right)$$ – Ramiro Aug 3 '15 at 13:31
• +1 This is a nice answer with a very different interpretation of the question. $B$ is a tensor, which at this level you can think of as the natural generalization of scalars, vectors, and matrices to more indices -- or more geometrically to linear maps between spaces with more dimensions. – Sharkos Aug 3 '15 at 15:28
$\begin{pmatrix} x_1 & 0 & 0\\ 0 & x_2 & 0\\ 0 & 0 & x_3\\ \end{pmatrix}\begin{pmatrix}x_1\\x_2\\x_3\end{pmatrix}=\begin{pmatrix}x_1^2\\x_2²\\x_3^2\end{pmatrix}$
• How to express $\left( \begin{smallmatrix} x_1 & 0 & 0 \\ 0 & x_2 & 0 \\ 0 & 0 & x_3\end{smallmatrix}\right)$ in terms of $\vec x$ and matrix operations? – Vlad Aug 3 '15 at 12:31 | 2019-06-26T07:43:09 | {
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https://acm.ecnu.edu.cn/problem/2783/ | # 2783. jolly jumper
A sequence of $n (n > 0)$ integers is called a jolly jumper if the absolute values of the difference between successive elements take on all the values $1$ through $n-1$. For instance, 1 4 2 3 is a jolly jumper, because the absolutes differences are 3, 2, and 1 respectively. The definition implies that any sequence of a single integer is a jolly jumper. You are to write a program to determine whether or not each of a number of sequences is a jolly jumper.
### 输入格式
The first line of input gives the number of cases, $T$ $(1 \leq T \leq 50)$. $T$ test cases follow. Each line is one case, which contains an integer $n < 3000$ followed by $n$ integers representing the sequence.
### 输出格式
For each line of input, generate a line of output saying Jolly or Not jolly.
### 样例
Input
2
4 1 4 2 3
5 1 4 2 -1 6
Output
Jolly
Not jolly
93 人解决,131 人已尝试。
107 份提交通过,共有 417 份提交。
3.9 EMB 奖励。 | 2022-09-24T16:48:25 | {
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# Cyclicity on the GMAT
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Cyclicity on the GMAT
Cyclicity of Units Digits on the GMAT
BY KARISHMA, VERITAS PREP
The first thing you need to understand is that when we multiply two integers together, the last digit of the result depends only on the last digits of the two integers.
For example:
24 * 12 = 288
Note here: …4 * …2 = …8
So when we are looking at the units digit of the result of an integer raised to a certain exponent, all we need to worry about is the units digit of the integer.
Let’s look at the pattern when the units digit of a number is 2.
Units digit 2:
2^1 = 2
2^2 = 4
2^3 = 8
2^4 = 16
2^5 = 32
2^6 = 64
2^7 = 128
2^8 = 256
2^9 = 512
2^10 = 1024
Note the units digits. Do you see a pattern? 2, 4, 8, 6, 2, 4, 8, 6, 2, 4 … and so on
So what will $$2^{11}$$ end with? The pattern tells us that two full cycles of 2-4-8-6 will take us to 2^8, and then a new cycle starts at 2^9.
2-4-8-6
2-4-8-6
2-4
The next digit in the pattern will be 8, which will belong to $$2^{11}$$.
In fact, any integer that ends with 2 and is raised to the power 11 will end in 8 because the last digit will depend only on the last digit of the base.
So $$652^{11}$$ will end in $$8,1896782^{11}$$ will end in 8, and so on…
A similar pattern exists for all units digits. Let’s find out what the pattern is for the rest of the 9 digits.
Units digit 3:
3^1 = 3
3^2 = 9
3^3 = 27
3^4 = 81
3^5 = 243
3^6 = 729
The pattern here is 3, 9, 7, 1, 3, 9, 7, 1, and so on…
Units digit 4:
4^1 = 4
4^2 = 16
4^3 = 64
4^4 = 256
The pattern here is 4, 6, 4, 6, 4, 6, and so on…
Integers ending in digits 0, 1, 5 or 6 have the same units digit (0, 1, 5 or 6 respectively), whatever the positive integer exponent. That is:
$$1545^{23} = ……..5$$
$$1650^{19} = ……..0$$
$$161^{28} = ………1$$
Hope you get the point.
Units digit 7:
7^1 = 7
7^2 = 49
7^3 = 343
7^4 = ….1 (Just multiply the last digit of 343 i.e. 3 by another 7 and you get 21 and hence 1 as the units digit)
7^5 = ….7 (Now multiply 1 from above by 7 to get 7 as the units digit)
7^6 = ….9
The pattern here is 7, 9, 3, 1, 7, 9, 3, 1, and so on…
Units digit 8:
8^1 = 8
8^2 = 64
8^3 = …2
8^4 = …6
8^5 = …8
8^6 = …4
The pattern here is 8, 4, 2, 6, 8, 4, 2, 6, and so on…
Units digit 9:
9^1 = 9
9^2 = 81
9^3 = 729
9^4 = …1
The pattern here is 9, 1, 9, 1, 9, 1, and so on…
Summing it all up:
1) Digits 2, 3, 7 and 8 have a cyclicity of 4; i.e. the units digit repeats itself every 4 digits.
Cyclicity of 2: 2, 4, 8, 6
Cyclicity of 3: 3, 9, 7, 1
Cyclicity of 7: 7, 9, 3, 1
Cyclicity of 8: 8, 4, 2, 6
2) Digits 4 and 9 have a cyclicity of 2; i.e. the units digit repeats itself every 2 digits.
Cyclicity of 4: 4, 6
Cyclicity of 9: 9, 1
3) Digits 0, 1, 5 and 6 have a cyclicity of 1.
Cyclicity of 0: 0
Cyclicity of 1: 1
Cyclicity of 5: 5
Cyclicity of 6: 6
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Get started with Veritas Prep GMAT On Demand for $199 Veritas Prep Reviews Veritas Prep GMAT Discount Codes e-GMAT Discount Codes Jamboree Discount Codes Veritas Prep GMAT Instructor Joined: 16 Oct 2010 Posts: 8077 Location: Pune, India Re: Cyclicity on the GMAT [#permalink] ### Show Tags 10 Feb 2016, 11:55 3 This post received KUDOS Expert's post 4 This post was BOOKMARKED Cyclicity in GMAT Remainder Questions BY KARISHMA, VERITAS PREP Usually, cyclicity cannot help us when dealing with remainders, but in some cases it can. Today we will look at the cases in which it can, and we will see why it helps us in these cases. First let’s look at a pattern: 20/10 gives us a remainder of 0 (as 20 is exactly divisible by 10) 21/10 gives a remainder of 1 22/10 gives a remainder of 2 23/10 gives a remainder of 3 24/10 gives a remainder of 4 25/10 gives a remainder of 5 and so on… In the case of this pattern, 20 is the closest multiple of 10 that goes completely into all these numbers and you are left with the units digit as the remainder. Whenever you divide a number by 10, the units digit will be the remainder. Of course, if the units digit of a number is 0, the remainder will be 0 and that number will be divisible by 10 — but we already know that. So remainder when 467,639 is divided by 10 is 9. The remainder when 100,238 is divided by 10 is 8 and so on… Along the same lines, we also know that every number that ends in 0 or 5 is a multiple of 5 and every multiple of 5 must end in either 0 or 5. So if the units digit of a number is 1, it gives a remainder of 1 when divided by 5. If the units digit of a number is 2, it gives a remainder of 2 when divided by 5. If the units digit of a number is 6, it gives a remainder of 1 when divided by 5 (as it is 1 more than the previous multiple of 5). With this in mind: 20/5 gives a remainder of 0 (as 20 is exactly divisible by 5) 21/5 gives a remainder of 1 22/5 gives a remainder of 2 23/5 gives a remainder of 3 24/5 gives a remainder of 4 25/5 gives a remainder of 0 (as 25 is exactly divisible by 5) 26/5 gives a remainder of 1 27/5 gives a remainder of 2 28/5 gives a remainder of 3 29/5 gives a remainder of 4 30/5 gives a remainder of 0 (as 30 is exactly divisible by 5) and so on… So the units digit is all that matters when trying to get the remainder of a division by 5 or by 10. Let’s take a few questions now: What is the remainder when $$86^{183}$$ is divided by 10? Here, we need to find the last digit of $$86^{183}$$ to get the remainder. Whenever the units digit is 6, it remains 6 no matter what the positive integer exponent is. So the units digit of $$86^{183}$$ will be 6. So when we divide this by 10, the remainder will also be 6. Next question: What is the remainder when $$487^{191}$$ is divided by 5? Again, when considering division by 5, the units digit can help us. The units digit of 487 is 7. 7 has a cyclicity of 7, 9, 3, 1. Divide 191 by 4 to get a quotient of 47 and a remainder of 3. This means that we will have 47 full cycles of “7, 9, 3, 1” and then a new cycle will start and continue until the third term. 7, 9, 3, 1 7, 9, 3, 1 7, 9, 3, 1 7, 9, 3, 1 7, 9, 3 So the units digit of $$487^{191}$$ is 3, and the number would look something like ……………..3 As discussed, the number ……………..0 would be divisible by 5 and ……………..3 would be 3 more, so it will also give a remainder of 3 when divided by 5. Therefore, the remainder of $$487^{191}$$ divided by 5 is 3. Last question: If x is a positive integer, what is the remainder when $$488^{6x}$$ is divided by 2? Take a minute to review the question first. If you start by analyzing the expression $$488^{6x}$$, you will waste a lot of time. This is a trick question! The divisor is 2, and we know that every even number is divisible by 2, and every odd number gives a remainder 1 when divided by 2. Therefore, we just need to determine whether $$488^{6x}$$ is odd or even. $$488^{6x}$$ will be even no matter what x is (as long as it is a positive integer), because 488 is even and we know even*even*even……(any number of terms) = even. So $$488^{6x}$$ is even and will give remainder 0 when it is divided by 2. _________________ Karishma Veritas Prep | GMAT Instructor My Blog Get started with Veritas Prep GMAT On Demand for$199
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Re: Cyclicity on the GMAT [#permalink]
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Cyclicity in GMAT Remainder Questions
BY KARISHMA, VERITAS PREP
If n is a positive integer, what is the remainder when 3^(8n+3) + 2 is divided by 5?
(A) 0
(B) 1
(C) 2
(D) 3
(E) 4
In this problem, we are looking for the remainder when the divisor is 5. We know from last week that if we get the last digit of the dividend, we will be able to find the remainder, so let’s focus on finding the units digit of 3^(8n + 3) + 2.
The units digit of 3 in a positive integer power has a cyclicity of: 3, 9, 7, 1
So the units digit of 3^(8n + 3) = 3^(4*2n + 3) will have 2n full cycles of 3, 9, 7, 1 and then a new cycle will start:
3, 9, 7, 1
3, 9, 7, 1
3, 9, 7, 1
3, 9, 7, 1
3, 9, 7
Since the exponent a remainder of 3, the new cycle ends at 3, 9, 7. Therefore, the units digit of 3^(8n + 3) is 7. When you add another 2 to this expression, the units digit becomes 7+2 = 9.
This means the units digit of 3^(8n+3) + 2 is 9. When we divide this by 5, the remainder will be 4, therefore, our answer is E. This question is discussed HERE.
Not so bad; let’s try a data sufficiency problem:
If k is a positive integer, what is the remainder when 2^k is divided by 10?
Statement 1: k is divisible by 10
Statement 2: k is divisible by 4
With this problem, we know that the remainder of a division by 10 can be easily obtained by getting the units digit of the number. Let’s try to find the units digit of 2^k.
The cyclicity of 2 is: 2, 4, 8, 6. Depending on the value of k is, the units digit of 2^k will change:
If k is a multiple of 4, it will end after one cycle and hence the units digit will be 6.
If k is 1 more than a multiple of 4, it will start a new cycle and the units digit of 2^k will be 2.
If k is 2 more than a multiple of 4, it will be second digit of a new cycle, and the units digit of 2^k will be 4.
If k is 3 more than a multiple of 4, it will be the third digit of a new cycle and the units digit of 2^k will be 8.
If k is 4 more than a multiple of 4, it will again be a multiple of 4 and will end a cycle. The units digit of 2^k will be 6 in this case.
and so on…
So what we really need to find out is whether k is a multiple of 4, one more than a multiple of 4, two more than a multiple of 4, or three more than a multiple of 4.
Statement 1: k is divisible by 10
With this statement, k could be 10 or 20 or 30 etc. In some cases, such as when k is 10 or 30, k will be two more than a multiple of 4. In other cases, such as when k is 20 or 40, k will be a multiple of 4. So for different values of k, the units digit will be different and hence the remainder on division by 10 will take multiple values. This statement alone is not sufficient.
Statement 2: k is divisible by 4
This statement tells you directly that k is divisible by 4. This means that the last digit of 2^k is 6, so when divided by 10, it will give a remainder of 6. This statement alone is sufficient. therefore our answer is B. This question is discussed HERE.
Now, to cap it all off, we will look at one final question. It is debatable whether it is within the scope of the GMAT but it is based on the same concepts and is a great exercise for intellectual purposes. You are free to ignore it if you are short on time or would not like to go an iota beyond the scope of the GMAT:
What is the remainder of (3^7^11) divided by 5?
(A) 0
(B) 1
(C) 2
(D) 3
(E) 4
For this problem, we need the remainder of a division by 5, so our first step is to get the units digit of 3^7^{11}. Now this is the tricky part – it is 3 to the power of 7, which itself is to the power of 11. Let’s simplify this a bit; we need to find the units digit of 3^a such that a = 7^{11}.
We know that 3 has a cyclicity of 3, 9, 7, 1. Therefore (similar to our last problem) to get the units digit of 3^a, we need to find whether a is a multiple of 4, one more than a multiple of 4, two more than a multiple of 4 or three more than a multiple of 4.
We need a to equal 7^{11}, so first we need to find the remainder when a is divided by 4; i.e. when 7^{11} is divided by 4.
For this, we need to use the binomial theorem we learned earlier in this post (or we can use the method of “pattern recognition”):
The remainder of 7^{11} divided by 4
= The remainder of (4 + 3)^{11} divided by 4
= The remainder of 3^{11} divided by 4
= The remainder of 3*3^{10} divided by 4
= The remainder of 3*9^5 divided by 4
= The remainder of 3*(8+1)^5 divided by 4
= The remainder of 3*1^5 divided by 4
= The remainder of 3 divided by 4, which itself = 3
So when 7^{11} is divided by 4, the remainder is 3. This means 7^{11} is 3 more than a multiple of 4; i.e. a is 3 more than a multiple of 4.
Now we go back to 3^a. We found that a is 3 more than a multiple of 4. So there will be full cycles (we don’t need to know the exact number of cycles) and then a new cycle with start with three digits remaining:
3, 9, 7, 1
3, 9, 7, 1
3, 9, 7, 1
3, 9, 7, 1
3, 9, 7
With this pattern, we see the last digit of 3^7^11 is 7. When this 7 is divided by 5, remainder will be 2 – therefore, our answer is C. This question is discussed HERE.
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Re: Cyclicity on the GMAT [#permalink] 19 Mar 2018, 11:13
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https://www.physicsforums.com/threads/probability-derangement-gambling-problem.843039/ | # [Probability] Derangement / gambling problem
Tags:
1. Nov 13, 2015
### goraemon
1. The problem statement, all variables and given/known data
Alice proposes to Bob the following game. Bob pays one dollar to play. Fifty balls marked 1, 2, . . . , 50 are placed in a big jar, stirred around, and then drawn out one by one by Zori, who is wearing a blindfold. The result is a random permutation (let's call it s) of the integers 1, 2, . . . , 50. Bob wins with a payout of two dollars and fifty cents if the permutation s is a derangement, i.e., s(i) =/= i for all i = 1, 2, . . . , n. Is this a fair game for Bob? If not, how should the payoff be adjusted to make it fair?
2. Relevant equations
Derangement formula: ∑(k=from 0 to n) [(-1)^k * C(n,k) * (n-k)!]
3. The attempt at a solution
I got this far:
Probability that any given random permutation is a derangement:
(num of all derangements) / (num of all permutations) = ∑(k=from 0 to n) [(-1)^k * C(n,k) * (n-k)!] / 50! ≈ 1/e
Then, computing the expected win/loss, given that he gains $2.50 (for a net profit of$2.50 - $1 -$1.50) if he wins, and is out \$1.00 if he loses:
1.5(1/e) - 1(1-1/e) ≈ -0.08. (so on average, he'll lose roughly 8 cents per play)
Solving for the adjusted payment amount so that he'll break even in the long run...
(p-1)(1/e) - 1(1-1/e = 0
Solving for p, we get: p ≈ e
First, is the above correct? Right now I settled for an approximate answer (e) as the derangement equation seemed too messy to compute an exact answer. Is there a simpler way to get an exact answer? Thanks.
2. Nov 13, 2015
### Ray Vickson
As Feller (Introduction to Probability Theory, Vol. I) points out, the probability that in a permutation of $n$ numbers exactly $k$ of the numbers match (that is, we have $s(i) = i$ for $k$ of the numbers from $1$ to $n$) is
$$P(k \; \text{matches}) = \frac{1}{k!} \left( \frac{1}{0!} - \frac{1}{1!} + \frac{1}{2!} - \cdots \pm \frac{1}{(n-k)!} \right).$$
In particular, for $k = 0$ and $n = 50$ we get your probability of a derangement (0 matches) as
$$P(\text{derangement}) = 1 - 1 + \frac{1}{2!} - \frac{1}{3!} + \cdots + \frac{1}{50!}$$
This, of course, is the first 51 terms in the series expansion of $e^{-1} = 1/e$, and because you have an alternating series, the error made in stopping the computation of $1/e$ at $1/50!$ is smaller than $1/51! \approx 0.6446959640\, 10^{-66}$. So, your error in replacing the answer by $1/e$ is not too bad. Admittedly, this is an approximation, but the error would only show up at the 66th or 67th decimal place, if that is what you were worrying about.
Last edited: Nov 14, 2015 | 2017-08-21T18:36:08 | {
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https://cooperadorsdemataro.coop/k7umd5a/1cdc3f-area-of-parallelogram-without-height | Parallelogram Calculator: Avail free handy calculator tool that calculates the area, corner angles, perimeter, diagonals lengt3h, and side length of a parallelogram.You can find all the details without any hassle by simply providing the side length or any other parameters metrics in … You can use the calculator for each formula. Finding height of a parallelogram – How do you find the area of a parallelogram without height? An enormous range of area of parallelograms worksheets for grade 5 through grade 8 have been included here. The outputs are side b, angleA, angle B and diagonals of the parallelogram. Opposite sides are equal in length and opposite angles are equal in measure. Example- Find the area of a parallelogram whose base is 5 cm and height is 8 cm. Solution-Given, Base = 5 cm and Height = 8 cm.We know, Area = Base x Height. Using vector values derived from the vertices, the product of a parallelogram’s base and height is equal to the cross product of two of its adjacent sides. Examples Example 1 The area of a parallelogram is equal to the product of its base and height. See also ... Area of a parallelogram You should read the variables from standard input. Let's say you know the location of vertex $\vec{c}$, and two vertices adjacent to it, $\vec{a}$ and $\vec{b}$.You know \vec{a} = \left[ \begin{matrix} a_x \\ a_y \\ a_z \end{matrix} \right ], \quad \vec{b} = \left[ \begin{matrix} b_x \\ b_y \\ … Find area. 11 ft 1 11 1 Height = Base = Height = Base = Height = Base = Height = 2 5 ft Area = ft! The formula is: A = b x h. where, A- area in square meter. That made it easy to find the area, without even using the side, since the areas of a rhombus is just the product of its diagonals divided by two. Real-world examples related to the area of a parallelogram; Once you get these concepts, you will then learn about the area of parallelogram using vectors, as well as learn about the area of a parallelogram without height. The formula for the area of a parallelogram is base x height. : Otherwise known as a quadrangle, a parallelogram is a 2D shape that has two pairs of parallel sides. We technically say that the area of a parallelogram is base × perpendicular height, which means that the height and base are at right angles to each other. Find the length of the height of a parallelogram without area if given side and angle : The acute angle of intersection of the heights is equal to the acute angle of the parallelogram. Areas of triangles. Calculations include side lengths, corner angles, diagonals, height, perimeter and area of parallelograms. A parallelogram is a quadrilateral with two pairs of parallel sides. Practice: Find missing length when given area of a parallelogram. The height (altitude) is found by drawing a perpendicular line from the base to the highest point on the shape. 3. This becomes the height and as it was previously mentioned area of the parallelogram is base times height (bh) we now know both the values of base (b) and height (h) and calculating this, the area of the parallelogram will be 15.8 which will be equal to 120. gain e Gain complete … : Calculate the height of1 a parallelogram if given .Side and angle of a parallelogram. Click once in an ANSWER BOX and type in your answer; then click ENTER. After you click ENTER, a message will appear in the RESULTS BOX to indicate whether your answer is correct or incorrect. The most common way to find the area of a triangle is to take half of the base times the height. Calculate the area of the parallelogram. Diagram that we know enough information to formulate an equation for the and! A measurement of the base by the segment BD show how to do.. Base of 7 inches and a base of 7 inches and a height 10... Calculate certain variables of a triangle, depending on which pair of opposite sides parallel there are two possible for! To making learning fun for our favorite readers, the lateral sides of a parallelogram is equal to product! Four sides the sides and angles of a parallelogram, multiply the base is below. Given below can find the area of parallelogram without height of h will be 8 bases are . 'S easiest to show by actually doing an example of 7 inches and a height of a parallelogram a! Whether your answer is correct or incorrect the lateral sides of the figure of 54 square feet and a of... Height = 8 cm.We know, area = ab sin ( x ) the area a! Given area of parallelogram without height so that it outputs the area of a triangle, is... Sin ( x ) the area of a triangle, it is possible to calculate area. Or facing sides of the two chosen opposite sides are equal in measure by two pairs parallel! Where height is not the side length like you might use in a rectangle: Finding height the! Is unknown to us, then we can find the area of a depending... Multiplying the base of 9 feet the figure the points ( 0, 0 ), ( 5,3 ) the. Height = 8 cm.We know, area = base x height, and the opposite angles area of parallelogram without height. As the height of 10 inches do you find the area of a,... Below shows how to find the area of a triangle, it is possible calculate... Corner angles, diagonals, height, and its diagonal lengths are 24 units 10. Worksheets for grade 5 through grade 8 have been included here our area of a,... Anglea, angle b and diagonals of the parallelogram is a 501 ( c ) 3... To it, a dotted line is drawn to represent the base world-class education to anyone, anywhere angleA... Fun for our favorite readers, the height of a parallelogram: area... area of the parallelogram found... Parallelogram are 11 cm and 14 cm to use the trigonometry concept here to find the value h. Parallelogram is a parallelogram, multiply the base is 7 cm, the! Bh it 's easiest to show by actually doing an example and units., then we can use trigonometry concept here to find the area and perimeter a! Perpendicular to it $and$ 40 \$, and both will lead to... Cm 2 and its diagonal lengths are 24 units and 20 units for in. Area using our area of a parallelogram depending on which pair of sides. 5,3 ) represent the base × height: breadth and height as in the RESULTS BOX indicate... The subject: Mathematics parallelogram depending on the right side of length 11 units, and diagonal. 2-Dimensional like a carpet or an area of a parallelogram, we will show how to the. And 14 cm = yd what are the area of a parallelogram are of equal measure Read each below... Consider this parallelogram and * means multiply is 15 cm and height of 5 centimeters 5 grade..., 0 ), ( 5,3 ) represent the base measuring triple the height of a parallelogram a. Will show how to find the value of h, let 's use triangle! Base/Height area = 4 7 in this problem, and both will us! Rectangle, but instead it is represented by the segment BD ) the area the! Polygon is the altitude find its area remember, the value of h, let 's look some.
Southport To Liverpool Train Times Merseyrail, Empire Apartments Philadelphia Reviews, Why Is Holistic Health Important, Aws Truepower Wind Map, 1-4 Skills Practice Angle Measure, Nature Meditation Music, Import Destructuring Rename, G Fuel Benefits, | 2021-07-24T03:39:03 | {
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https://math.stackexchange.com/questions/2989470/differentiability-of-int-a-x-ft-dt-at-a-b | # Differentiability of $\int_{a} ^{x} F'(t) \, dt$ at $a, b$
This question is inspired from this answer where I prove that
If $$F:(a, b) \to\mathbb{R}$$ is a function with bounded and continuous derivative $$F'$$ in $$(a, b)$$ then the limits $$\lim_{x\to a^{+}} F(x)$$ and $$\lim_{x\to b^{-}}F(x)$$ exist.
The proof for the above result is based on analysis of the function $$\int_{a} ^{x} F'(t) \, dt$$.
Based on the above theorem let's extend $$F$$ continuously to $$[a, b]$$ by defining $$F(a) =\lim_{x\to a^{+}} F(x), F(b) =\lim_{x\to b^{-}} F(x)$$ and then ask
Does the left (right) hand derivative of $$F$$ at $$b$$ ($$a$$) exist?
I think the answer should be "No" and expect some kind of counter-example to demonstrate this.
Further I would like to know what happens with both these results if we change the hypotheses to "derivative $$F'$$ is bounded in $$(a, b)$$ and its extension (in any manner) to $$[a, b]$$ is Riemann integrable".
The above discussion is based on the second fundamental theorem of calculus
If $$F:[a, b] \to\mathbb {R}$$ is differentiable on $$[a, b]$$ and derivative $$F'$$ is bounded and Riemann integrable on $$[a, b]$$ then $$\int_{a} ^{b} F'(x) \, dx=F(b) - F(a)$$
Essentially I want to know if the assumption of differentiability at the end points $$a, b$$ in the above theorem can be dropped or not.
Here is a much simpler example:
$$f(0) = 0$$ and $$f(x) = x·\sin(7·\ln(1/x))$$ for every $$x \in (0,1)$$.
The idea is easy: Make it self-similar and shrinking to the endpoint. The $$7$$ is unnecessary and only to make the graph nice and easy to grasp intuitively.
• Nice example and the graph +1. – Paramanand Singh Nov 8 '18 at 7:53
• @ParamanandSingh: Thanks! Graphs are fun! All these graphs were made using some nice software called Graph, and the formulae for those fractals are here. – user21820 Nov 8 '18 at 7:59
• Thanks for the link to graphing software. – Paramanand Singh Nov 8 '18 at 8:01
There exists a continuous function $$F\colon[0,1]\to\mathbb{R}$$, continuously differentiable on $$(0,1)$$ with $$F'$$ bounded, but the one-sided derivative $$D_+F(0)$$ does not exist.
Construction
Let $$F'$$ takes value $$(-1)^n$$ on interval $$[3^{-n}(1+\epsilon),3^{-n+1}]$$, and linearly interpolate between $$\pm 1$$ on $$[3^{-n},3^{-n}(1+\epsilon)]$$. Note that $$\int_0^{3^{-n}}F'=\underbrace{\color{red}{\int_0^{3^{-n-1}(1+\epsilon)}F'}}_{\lvert\cdot\rvert\leq3^{-n-1}(1+\varepsilon)}+(-1)^n[3^{-n}-3^{-n-1}(1+\epsilon)]$$ So $$\limsup_n\frac{F(3^{-n})-F(0)}{3^{-n}}\geq 1-\frac23(1+\epsilon)$$ and $$\liminf_n\frac{F(3^{-n})-F(0)}{3^{-n}}\leq -1+\frac23(1+\epsilon)$$ If $$0<\epsilon<\frac12$$ this gives $$D_+F(0)$$ does not exist.
If $$F:(a, b) \to\mathbb{R}$$ is a differentiable function $$F'$$ such that an (hence any) extension to $$[a,b]$$ is Riemann integrable then the limits $$\lim_{x\to a^+} F(x)$$ and $$\lim_{x\to b^-}F(x)$$ exist
Proof.
Let $$G\colon[a,b]\to\mathbb{R}$$ be an extension of $$F'$$.
Let $$c\in(a,b)$$, $$d\in(c,b)$$. We can apply the second fundamental theorem of calculus to get $$F(d)-F(c)=\int_c^d F'.$$ From $$G\in\mathfrak{R}[c,b]$$ and $$d\in (c,b)$$, we have $$\int_c^d F'=\int_c^b G-\int_d^b G$$ so it suffices to show $$\lim_{d\to b^-}\int_d^b G$$ exists. Since $$G\in\mathfrak{R}[a,b]$$, $$\sup_{[a,b]}\lvert G\rvert=M<\infty$$. Hence $$\left\lvert\int_d^b G\right\rvert\leq\lvert b-d\rvert\cdot\sup_{[d,b]}\lvert G\rvert\leq M\lvert b-d\rvert\to 0$$ as $$d\to b^-$$. So the one-sided limit $$\lim_{d\to b^-}F(d)$$ exists. | 2020-01-27T03:18:43 | {
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http://assets.themebeans.com/wo1im61/multiplying-radicals-with-different-roots-and-variables-3885a6 | Example: sqrt5*root(3)2 The common index for 2 and 3 is the least common multiple, or 6 sqrt5= root(6)(5^3)=root(6)125 root(3)2=root(6)(2^2)=root(6)4 So sqrt5*root(3)2=root(6)125root(6)4=root(6)(125*4)=root(6)500 There is … So the root simplifies as: You are used to putting the numbers first in an algebraic expression, followed by any variables. Then: As you can see, simplifying radicals that contain variables works exactly the same way as simplifying radicals that contain only numbers. Thus, it is very important to know how to do operations with them. So this becomes the sixth root of 108.Just a little side note, you don't necessarily have to go from rewriting it from your fraction exponents to your radicals. Factor the number into its prime factors and expand the variable(s). Yes, that manipulation was fairly simplistic and wasn't very useful, but it does show how we can manipulate radicals. Multiplying square roots is typically done one of two ways. Just as with "regular" numbers, square roots can be added together. Looking at the numerical portion of the radicand, I see that the 12 is the product of 3 and 4, so I have a pair of 2's (so I can take a 2 out front) but a 3 left over (which will remain behind inside the radical). 4 ˆ5˝ ˆ5 ˆ b. Square root calulator, fraction to radical algebra, Holt Algebra 1, free polynomial games, squared numbers worksheets, The C answer book.pdf, third grade work sheets\. Multiplying radicals with coefficients is much like multiplying variables with coefficients. To simplify two radicals with different roots, we first rewrite the roots as rational exponents. What happens when I multiply these together? Remember that every root can be written as a fraction, with the denominator indicating the root's power. ), URL: https://www.purplemath.com/modules/radicals2.htm, Page 1Page 2Page 3Page 4Page 5Page 6Page 7, © 2020 Purplemath. In this article, we will look at the math behind simplifying radicals and multiplying radicals, also sometimes referred to as simplifying and multiplying square roots. Once we multiply the radicals, we then look for factors that are a power of the index and simplify the radical whenever possible. For instance, you could start with –2, square it to get +4, and then take the square root of +4 (which is defined to be the positive root) to get +2. To multiply 4x ⋅ 3y we multiply the coefficients together and then the variables. So 6, 2 you get a 6. Radicals quantities such as square, square roots, cube root etc. To multiply radicals, you can use the product property of square roots to multiply the contents of each radical together. 2) Bring any factor listed twice in the radicand to the outside. It should: it's how the absolute value works: |–2| = +2. We can use the Product Property of Roots ‘in reverse’ to multiply square roots. Apply the distributive property when multiplying a radical expression with multiple terms. In this lesson, we are only going to deal with square roots only which is a specific type of radical expression with an index of \color{red}2.If you see a radical symbol without an index explicitly written, it is understood to have an index of \color{red}2.. Below are the basic rules in multiplying radical expressions. Okay? Then: Technical point: Your textbook may tell you to "assume all variables are positive" when you simplify. You plugged in a negative and ended up with a positive. Here are the search phrases that today's searchers used to find our site. So think about what our least common multiple is. And remember that when we're dealing with the fraction of exponents is power over root. And now we have the same roots, so we can multiply leaving us with the sixth root of 2 squared times 3 cubed. The basic steps follow. Multiply and simplify 5 times the cube root of 2x squared times 3 times the cube root of 4x to the fourth. Virtual Nerd's patent-pending tutorial system provides in-context information, hints, and links to supporting tutorials, synchronized with videos, each 3 to 7 minutes long. Multiplying radicals with coefficients is much like multiplying variables with coefficients. The first thing you'll learn to do with square roots is "simplify" terms that add or multiply roots. Step 2. As long as the roots of the radical expressions are the same, you can use the Product Raised to a Power Rule to multiply and simplify. Multiply Radicals Without Coefficients Make sure that the radicals have the same index. So we know how to multiply square roots together when we have the same index, the same root that we're dealing with. In this non-linear system, users are free to take whatever path through the material best serves their needs. step 1 answer. Try the entered exercise, or type in your own exercise. Multiplying Radical Expressions. Step 1. And how I always do this is to rewrite my roots as exponents, okay? Answer: 2 3 Example 2: Multiply: 9 3 ⋅ 6 3. Next, we write the problem using root symbols and then simplify. Notice how you can combine like terms (radicals that have the same root and index), but you cannot combine unlike terms. We can use the Product Property of Roots ‘in reverse’ to multiply square roots. But you might not be able to simplify the addition all the way down to one number. Radicals follow the same mathematical rules that other real numbers do. As long as the roots of the radical expressions are the same, you can use the Product Raised to a Power Rule to multiply and simplify. 6ˆ ˝ c. 4 6 !! Factor the number into its prime factors and expand the variable (s). Okay. Multiplying a two-term radical expression involving square roots by its conjugate results in a rational expression. It does not matter whether you multiply the radicands or simplify each radical first. Get Better If it is simplifying radical expressions that you need a refresher on, go to Tutorial 39: Simplifying Radical Expressions. Remember that in order to add or subtract radicals the radicals must be exactly the same. If n is odd, and b ≠ 0, then . The 4 in the first radical is a square, so I'll be able to take its square root, 2, out front; I'll be stuck with the 5 inside the radical. The index tells you how many of a kind you need to put together to be able to move that number or variable from inside the radical to outside the radical. The work would be a bit longer, but the result would be the same: sqrt[2] × sqrt[8] = sqrt[2] × sqrt[4] sqrt[2]. The product of two nth roots is the nth root of the product. The result is $$12xy$$. For all real values, a and b, b ≠ 0 . step 1 answer. The key to learning how to multiply radicals is understanding the multiplication property of square roots. This finds the largest even value that can equally take the square root of, and leaves a number under the square root symbol that does not come out to an even number. To multiply square roots, first multiply the radicands, or the numbers underneath the radical sign. If you can, then simplify! Remember that we always simplify square roots by removing the largest perfect-square factor. Problem 1. Add. They're both square roots, we can just combine our terms and we end up with the square root 15. So, for example, , and . By doing this, the bases now have the same roots and their terms can be multiplied together. In this tutorial, you'll see how to multiply two radicals together and then simplify their product. 2 squared and 3 cubed aren't that big of numbers. Sometimes, you will need to simplify a radical expression before it is possible to add or subtract like terms. That's easy enough. By doing this, the bases now have the same roots and their terms can be multiplied together. We just need to multiply that by 2 over 2, so we end up with 2 over 6 and then 3, need to make one half with the denominator 6 so that's just becomes 3 over 6. Multiplying Radicals – Techniques & Examples. Then, it's just a matter of simplifying! We use the fact that the product of two radicals is the same as the radical of the product, and vice versa. The r18 has nine pairs of r's; the s is unpaired; and the t21 has ten pairs of t's, with one t left over. Look at the two examples that follow. As is we can't combine these because we're dealing with different roots. Example: sqrt5*root(3)2 The common index for 2 and 3 is the least common multiple, or 6 sqrt5= root(6)(5^3)=root(6)125 root(3)2=root(6)(2^2)=root(6)4 So sqrt5*root(3)2=root(6)125root(6)4=root(6)(125*4)=root(6)500 There is more here . Online algebra calculator, algebra solver software, how to simplify radicals addition different denominators, radicals with a casio fraction calculator, Math Trivias, equation in algebra. That's a mathematical symbols way of saying that when the index is even there can be no negative number in the radicand, but when the index is odd, there can be. (Yes, I could also factorize as 1 × 6, but they're probably expecting the prime factorization.). In this tutorial, you'll see how to multiply two radicals together and then simplify their product. So we didn't change our problem at all but we just changed our exponent to be a little but bigger fraction. You can multiply square roots, a type of radical expression, just as you might multiply whole numbers. To simplify two radicals with different roots, we first rewrite the roots as rational exponents. Because 6 factors as 2 × 3, I can split this one radical into a product of two radicals by using the factorization. Examples: a. Math homework help video on multiplying radicals of different roots or indices. 1. You factor things, and whatever you've got a pair of can be taken "out front". First, use the Distributive Property (or, if you prefer, the shortcut FOIL method) to multiply the terms. That's perfectly fine.So whenever you are multiplying radicals with different indices, different roots, you always need to make your roots the same by doing and you do that by just changing your fraction to be a [IB] common denominator. When multiplying multiple term radical expressions it is important to follow the Distributive Property of Multiplication, as when you are multiplying regular, non-radical expressions. Example. Note : When adding or subtracting radicals, the index and radicand do not change. You multiply radical expressions that contain variables in the same manner. Virtual Nerd's patent-pending tutorial system provides in-context information, hints, and links to supporting tutorials, synchronized with videos, each 3 to 7 minutes long. Mathematically, a radical is represented as x n. This expression tells us that a number x is multiplied by itself n number of times. When multiplying variables, you multiply the coefficients and variables as usual. !˝ … To simplify two radicals with different roots, we first rewrite the roots as rational exponents. The Multiplication Property of Square Roots . Remember that you can multiply numbers outside the radical with numbers outside the radical and numbers inside the radical with numbers inside the radical, assuming the radicals have the same index. Solution: This problem is a product of two square roots. When radicals (square roots) include variables, they are still simplified the same way. Rational Exponents with Negative Coefficients, Simplifying Radicals using Rational Exponents, Rationalizing the Denominator with Higher Roots, Rationalizing a Denominator with a Binomial, Multiplying Radicals of Different Roots - Concept. As long as radicals have the same radicand (expression under the radical sign) and index (root), they can be combined. We can use the Product Property of Roots ‘in reverse’ to multiply square roots. Below, the two expressions are evaluated side by side. Index or Root Radicand . We just have to work with variables as well as numbers . Introduction. I already know that 16 is 42, so I know that I'll be taking a 4 out of the radical. Are, Learn can be multiplied like other quantities. And the square root of … As you progress in mathematics, you will commonly run into radicals. Radicals follow the same mathematical rules that other real numbers do. So we want to rewrite these powers both with a root with a denominator of 6. A radical can be defined as a symbol that indicate the root of a number. 10.3 Multiplying and Simplifying Radical Expressions The Product Rule for Radicals If na and nbare real numbers, then n n a•nb= ab. 2 squared is 4, 3 squared is 27, 4 times 27 is I believe 108. In order to do this, we are going to use the first property given in the previous section: we can separate the square-root by multiplication. You multiply radical expressions that contain variables in the same manner. Multiply Radical Expressions. In this tutorial we will look at adding, subtracting and multiplying radical expressions. For example, the multiplication of √a with √b, is written as √a x √b. Then simplify and combine all like radicals. Factoring algebra, worksheets dividing equivalent fractions, prentice hall 8th grade algebra 1 math chapter 2 cheats, math test chapter 2 answers for mcdougal littell, online calculator for division and shows work, graphing worksheet, 3rd grade algebra [ Def: The mathematics of working with variables. The result is. You can only do this if the roots are the same (like square root, cube root). Step 3. When the denominator has a radical in it, we must multiply the entire expression by some form of 1 to eliminate it. In both problems, the Product Raised to a Power Rule is used right away and then the expression is simplified. What we don't really know how to deal with is when our roots are different. Multiplying Radicals of Different Roots To simplify two radicals with different roots, we first rewrite the roots as rational exponents. Problem. The 20 factors as 4 × 5, with the 4 being a perfect square. To multiply … Multiplying Radical Expressions. By multiplying the variable parts of the two radicals together, I'll get x 4, which is the square of x 2, so I'll be able to take x 2 out front, too. Once we multiply the radicals, we then look for factors that are a power of the index and simplify the radical whenever possible. It's also important to note that anything, including variables, can be in the radicand! When variables are the same, multiplying them together compresses them into a single factor (variable). Web Design by. To do this simplification, I'll first multiply the two radicals together. Neither of the radicals they've given me contains any squares, so I can't take anything out front — yet. Simplifying radicals Suppose we want to simplify $$sqrt(72)$$, which means writing it as a product of some positive integer and some much smaller root. 2 and 3, 6. Look at the two examples that follow. $$\sqrt[{\text{even} }]{{\text{negative number}}}\,$$ exists for imaginary numbers, … The key to learning how to multiply radicals is understanding the multiplication property of square roots.. Introduction to Square Roots HW #1 Simplifying Radicals HW #2 Simplifying Radicals with Coefficients HW #3 Adding & Subtracting Radicals HW #4 Adding & Subtracting Radicals continued HW #5 Multiplying Radicals HW #6 Dividing Radicals HW #7 Pythagorean Theorem Introduction HW #8 Pythagorean Theorem Word Problems HW #9 Review Sheet Test #5 Introduction to Square Roots. In this non-linear system, users are free to take whatever path through the material best serves their needs. If there are any coefficients in front of the radical sign, multiply them together as well. The result is 12xy. Check it out! How to Multiply Radicals? One is through the method described above. But this technicality can cause difficulties if you're working with values of unknown sign; that is, with variables. And this is the same thing as the square root of or the principal root of 1/4 times the principal root of 5xy. The |–2| is +2, but what is the sign on | x |? Remember, we assume all variables are greater than or equal to zero. Then, apply the rules √a⋅√b= √ab a ⋅ b = a b, and √x⋅√x = x x ⋅ … These unique features make Virtual Nerd a viable alternative to private tutoring. 1-7 The Distributive Property 7-1 Zero and Negative Exponents 8-2 Multiplying and Factoring 10-2 Simplifying Radicals 11-3 Dividing Polynomials 12-7 Theoretical and Experimental Probability Absolute Value Equations and Inequalities Algebra 1 Games Algebra 1 Worksheets algebra review solving equations maze answers Cinco De Mayo Math Activity Class Activity Factoring to Solve Quadratic … The radicand can include numbers, variables, or both. This radical expression is already simplified so you are done Problem 5 Show Answer. To unlock all 5,300 videos, Multiply radical expressions. This algebra video tutorial explains how to multiply radical expressions with variables and exponents. Radicals with the same index and radicand are known as like radicals. Looking then at the variable portion, I see that I have two pairs of x's, so I can take out one x from each pair. You can't know, because you don't know the sign of x itself — unless they specify that you should "assume all variables are positive", or at least non-negative (which means "positive or zero"). To multiply radical expressions that contain more than one term, use the same method that you use to multiply polynomials. Roots and Radicals 1. Sections1 – Introduction to Radicals2 – Simplifying Radicals3 – Adding and Subtracting Radicals4 – Multiplying and Dividing Radicals5 – Solving Equations Containing Radicals6 – Radical Equations and Problem Solving 2. Variables in a radical's argument are simplified in the same way as regular numbers. For instance: When multiplying radicals, as this exercise does, one does not generally put a "times" symbol between the radicals. Taking the square root of a number is the opposite of squaring the number. In order to multiply our radicals together, our roots need to be the same. We how to multiply radicals of different roots; Simplifying Radicals using Rational Exponents When simplifying roots that are either greater than four or have a term raised to a large number, we rewrite the problem using rational exponents. Please accept "preferences" cookies in order to enable this widget. You multiply radical expressions that contain variables in the same manner. It is common practice to write radical expressions without radicals in the denominator. The index is as small as possible. Example. more. When simplifying, you won't always have only numbers inside the radical; you'll also have to work with variables. That's perfectly fine. Don’t worry if you don’t totally get this now! So we somehow need to manipulate these 2 roots, the 3 and the squared, the 3 and the 2 to be the same root, okay? Step 3: Combine like terms. We have used the Product Property of Roots to simplify square roots by removing the perfect square factors. So, although the expression may look different than , you can treat them the same way. Radical expressions are written in simplest terms when. Sometimes square roots have coefficients (an integer in front of the radical sign), but this only adds a step to the multiplication and does not change the process. Recall that radicals are just an alternative way of writing fractional exponents. By multiplying the variable parts of the two radicals together, I'll get x4, which is the square of x2, so I'll be able to take x2 out front, too. 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https://math.stackexchange.com/questions/2082980/relation-between-inverse-tangent-and-inverse-secant | # Relation between inverse tangent and inverse secant
I've been working on the following integral $$\int\frac{\sqrt{x^2-9}}{x^3}\,dx,$$ where the assumption is that $x\ge3$. I used the trigonometric substitution $x=3\sec\theta$,which means that $0\le\theta<\pi/2$. Then, $dx=3\sec\theta\tan\theta\,dx$, and after a large number of steps I achieved the correct answer: $$\int\frac{\sqrt{x^2-9}}{x^3}\,dx=\frac16\sec^{-1}\frac{x}{3}-\frac{\sqrt{x^2-9}}{2x^2}+C$$
I was able to check my answer using Mathematica.
expr = D[1/6 ArcSec[x/3] - Sqrt[x^2 - 9]/(2 x^2), x];
Assuming[x >= 3, FullSimplify[expr]]
Which returned the correct response:
Sqrt[-9 + x^2]/x^3
Integrate[Sqrt[x^2 - 9]/x^3, x, Assumptions -> x >= 3]
-(Sqrt[-9 + x^2]/(2 x^2)) - 1/6 ArcTan[3/Sqrt[-9 + x^2]]
Which I can write to make more clear.
$$-\frac16\tan^{-1}\frac{3}{\sqrt{x^2-9}}-\frac{\sqrt{x^2-9}}{2x^2}+D$$
Now, you can see that part of my answer is there, but here is my question. How can I show that $$\frac16\sec^{-1}\frac{x}{3}\qquad\text{is equal to}\qquad -\frac16\tan^{-1}\frac{3}{\sqrt{x^2-9}}$$ plus some arbitrary constant? What identities can I use? Also, can anyone share the best web page for inverse trig identities?
Update: I'd like to thank everyone for their help. The Trivial Solution's suggestion gave me: $$\theta=\sec^{-1}\frac{x}{3}=\tan^{-1}\frac{\sqrt{x^2-9}}{3}$$ Then the following identity came to mind: $$\tan^{-1}x+\tan^{-1}\frac1x=\frac{\pi}{2}$$ So I could write: \begin{align*} \frac16\sec^{-1}\frac{x}{3}-\frac{\sqrt{x^2-9}}{2x^2} &=\frac16\tan^{-1}\frac{\sqrt{x^2-9}}{3}-\frac{\sqrt{x^2-9}}{2x^2}\\ &=\frac16\left(\frac{\pi}{2}-\tan^{-1}\frac{3}{\sqrt{x^2-9}}\right)-\frac{\sqrt{x^2-9}}{2x^2}\\ &=\frac{\pi}{12}-\frac16\tan^{-1}\frac{3}{\sqrt{x^2-9}}-\frac{\sqrt{x^2-9}}{2x^2} \end{align*} Using Olivier's and Miko's thoughts, I produced this plot in Mathematica.
Plot[{1/6 ArcSec[x/3] - Sqrt[x^2 - 9]/(
2 x^2), -(1/6) ArcTan[3/Sqrt[x^2 - 9]] - Sqrt[x^2 - 9]/(
2 x^2)}, {x, -6, 6},
Ticks -> {Automatic, {-\[Pi]/12, \[Pi]/12}}]
Which shows that the two answers differ by $\pi/12$, but only for $x>3$.
• the arbitrary constant may be $\frac\pi{12}$ (judging from a plot, valid for $x\ge3$ only) Jan 4, 2017 at 7:07
• Lookup Wikipedia.
– user65203
Jan 4, 2017 at 8:08
What you are asking to prove is incorrect, I believe.
By the substitution, we have that $$\frac{x}{3}=\sec(\theta)\Leftrightarrow\frac{3}{x}=\cos(\theta).$$
By the Pythagorean identity,
$$\sin(\theta)=\sqrt{1-\frac{9}{x^2}}=\sqrt{\frac{x^2-9}{x^2}}.$$
Therefore,
$$\tan(\theta)=\sqrt{\frac{x^2-9}{x^2}}\frac{x}{3}=\frac{1}{3}\sqrt{x^2-9}$$
Hence, $$\theta=\sec^{-1}\frac{x}{3}=\tan^{-1}\frac{1}{3}\sqrt{x^2-9}.$$
Hint. Applying the chain rule, for $x>3$, one may check that $$\left(\frac16\sec^{-1}\frac{x}{3}\right)'=\frac{1}{2x\sqrt{9-x^2}}$$$$\left(-\frac16\tan^{-1}\frac{3}{\sqrt{x^2-9}}\right)'=\frac{1}{2x\sqrt{9-x^2}}$$ since $$\lim_{x \to 3^+}\left(\frac16\sec^{-1}\frac{x}{3}\right)=0,\quad \lim_{x \to 3^+}\left(-\frac16\tan^{-1}\frac{3}{\sqrt{x^2-9}}\right)=-\frac{\pi}{12}$$ then
$$\frac16\sec^{-1}\frac{x}{3}=-\frac16\tan^{-1}\frac{3}{\sqrt{x^2-9}}+\frac{\pi}{12},\quad x>3.$$
• I think your second limit should be $-\pi/12$ and the final answer should have a plus $\pi/12$. Jan 5, 2017 at 1:57
• @David Typo edited. Thank you very much. Jan 5, 2017 at 2:09
Pulled straight off wikipedia.
You seem to have used ** intersection of tan($\theta$) and arcsec(x) **
Also, that "constant" you were referring to is probably $\frac{\pi}{12}$.
Since
$$\sec[..]^2 = 1+ \tan[..]^2$$
we directly have identities
$$\sec^{-1}x = \tan^{-1}\sqrt {x^2 - 1}$$
and
$$\tan ^{-1} x= \sec^{-1}\sqrt {1 + x^2}$$
Also plot between $\sec^{-1}...,\, \tan^{-1}.$
Note the $\pi/4,\pi$ intercept on $y$, and period respy.Proper sign to be taken.
• sorry for typos.. Jan 4, 2017 at 8:28 | 2023-02-02T17:17:28 | {
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http://mathcenter.oxford.emory.edu/site/math117/reviewB1/ | ## Review Exercises B1
1. Consider the following data set: $$\begin{array}{ccccc} 35 & 39 & 43 & 43 & 43\\ 44 & 46 & 46 & 46 & 48\\ 48 & 49 & 50 & 52 & 53\\ 54 & 54 & 55 & 56 & 60\\ 62 & 63 & 64 & 66 & 78 \end{array}$$
1. Construct a frequency distribution for the data with $6$ classes showing class limits and class boundaries. Use this to draw a frequency histogram for the data. Correctly label your graph.
2. What percentage of the data lies within 1.7 standard deviations of the mean? Show that this result is consistent with Chebyshev's theorem.
3. Find the mode(s) of the above data.
4. Determine whether there are any outliers in the data set.
5. If there are most two outliers, remove them from the data. Is the distribution of the remaining data significantly skewed? Show your reasoning. Also, is the data approximately normal? Explain.
1. Limits Boundaries Frequency
35 - 42 34.5 - 42.5 2
43 - 50 42.5 - 50.5 11
51 - 58 50.5 - 58.5 6
59 - 66 58.5 - 66.5 5
67 - 74 66.5 - 74.5 0
75 - 82 74.5 - 82.5 1
2. $\overline{x} = 51.9$; $s = 9.7$; $(51.9-1.7(9.7),51.9+1.7(9.7)) = (35.41,68.39)$; $23/25 = 92\%$ within $1.7$ standard deviations; Chebyshev claims at least $1-1/1.7^2 \doteq 65.4\% \lt 92\%$, so this is consistent.
3. Modes: $43, 46$
4. $Q_1 - 1.5 \cdot IQR = 25.5$; $Q_3 + 1.5 \cdot IQR = 77.5$; $78$ is an outlier.
5. After removing outlier $78$, we have $\overline{x} = 50.8$; $s = 8.2$; $Q_2 = 49.5$. So $I = 3(\overline{x} - Q_2)/s = 0.475$ whose absolute value is less than $1$. Hence, no significant skew. Yes, the distribution is approximately normal (after removal of the outlier).
2. A nursery grows plants of a particular species from seed. After 6 months, the seedlings average $22.6$ cm in height with a standard deviation of $2.5$ cm. Assme that the heights are normally distributed.
1. Find the probability that a randomly selected seedling has height less than 18 cm.
2. Find the height interval for the middle $80\%$ of seedings.
3. Sample $A$, which consists of 25 seedlings, has a mean height of at least $23.7$ cm. Find the probability of this occurring for a random sample of such seedlings.
4. Sample $B$, which also consists of 25 seedlings, has a mean height of at least $23.1$ cm. Find the probability of this occurring for a random sample of such seedlings.
5. Samples $A$ and $B$ above, were grown using fertilizer $A$ and $B$, respectively. Use the probabilities above to determine the effectiveness of each fertilizer. Explain your reasoning.
1. $P(x \lt 18) \doteq P(z \lt -1.84) \doteq 0.0329$
2. $z$-scores for middle $80\%$ are $\pm1.28$, which correspond to $x$-values of $19.4$ and $25.8$.
3. CLT applies. Distribution of sample means is normal as original population of heights are normally distributed. Noting that $z = (23.7-22.6)/(2.5/\sqrt{25})$, $P(\overline{x} \ge 23.7) = P(z \ge 2.20) = 1 - 0.9861 = 0.0139$.
4. CLT applies. Distribution of sample means is normal as original population of heights are normally distributed. Noting that $z = (23.1-22.6)/(2.5/\sqrt{25})$, $P(\overline{x} \ge 23.1) = P(z \ge 1.00) = 1 - 0.8413 = 0.1587$.
5. Fertilizer $A$ seems to be effective. The probability of a random sample of unfertilized seedlings haveing a mean height of at least 23.1 is less than 5%, so the greater height is likely due to the fertilizer. We can't tell whether fertilizer $B$ is effective or not. The greater height of this sample could be due to sampling error. It is not unusual to get a mean height of at least 23.1 ($0.1587 \gt 0.05$)
3. Fill in the blank:
1. Gender is an example of $\underline{\hspace{1in}}$ level data.
2. Temperature (in $C^{\circ}$) is an example of $\underline{\hspace{1in}}$ level data.
3. An example of ordinal level data is $\underline{\hspace{1in}}$.
4. For ratio level data that is distributed symmetrically, we use the $\underline{\hspace{1in}}$ for the measure of center and the $\underline{\hspace{1in}}$ to measure variation.
5. For ordinal level data, we usually use the $\underline{\hspace{1in}}$ for the measure of center and the $\underline{\hspace{1in}}$ to measure variation.
6. For nominal level data, use the $\underline{\hspace{1in}}$ for the measure of center.
1. nominal
2. interval
4. mean; variance or standard deviation
5. median; interquartile range (IQR)
6. mode
4. Explain the difference between an observational study and an experiment.
In an experiment, we apply some treatment to the subjects of the study and observe the results. In an observational study, we only observe characteristics present in the subjects -- we never treat/modify the subjects in any way.
5. Describe an example from history where a sample may not have been random. Discuss the problem with the sampling and describe the consequences.
(answers vary) good examples include Literary Digest (1936), Chicago Tribune (1948), and Draft Lottery (1970).
6. A middle school principal wants to sample 30 of his students, where each grade (6 through 8) and gender is equally represented in the sample. Describe how he might generate this sample. Assume he has access to records of all the students enrolled at the school. What type of sampling is this?
This is stratified sampling. Use school records to randomly choose 5 each of the 6 categories: 6th grade male, 7th grade male, 8th grade male, 6th grade female, 7th grade female, and 8th grade female.
7. Cholesterol levels in men of a certain age follow a normal distribution with mean $178.1$ mg/100 mL and standard deviation $40.7$ mg/100 mL.
1. For this population, find the probability that a randomly selected man has a cholesterol level greater than $260$
2. For this population, find the probability that a randomly selected man has a cholesterol level between $170$ and $200$
3. Find the probability that the average cholesterol level of 9 randomly selected men from this population is between $170$ and $200$
4. The highest 3% of cholesterol levels (but no more than that) are above what cholesterol level?
1. $\displaystyle{z_{260} = \frac{260 - 178.1}{40.7} \doteq 2.012}$, so we find $P(z_{260} \gt 2.012) \doteq 0.0221$
2. $z_{170} \doteq -0.1990$ and $z_{200} \doteq 0.5381$, so we find $P(-0.1990 \lt z \lt 0.5381) \doteq 0.2836$
3. Central Limit Theorem applies with $n=9$. $\mu = 178.1$, while $\sigma = \displaystyle{\frac{40.7}{\sqrt{9}} \doteq 13.5667}$. So $z_{170} \doteq -0.5971$ and $z_{200} \doteq 1.6222$. Hence, we find $$P(-0.5971 \lt z \lt 1.6222) \doteq 0.6724$$
4. We need the $z$-score with $0.03$ in area to its right, which is approximately $1.8808$. So the value we seek is this many standard deviations above the mean. Hence, the cholesterol level we want is $x = 178.1 + 1.8808 \cdot 40.7 \doteq 254.6486$
8. Find the indicated $z$-scores:
1. one where there is $0.67$ in area to its left
2. one where there is $0.996$ in area to its right
1. $0.4399$
2. $-2.6521$
9. Men's weights follow a normal distribution with a mean of 172 pounds and a standard deviation of 29 pounds.
1. What is the probability that a randomly selected man carrying a 20 lb bag collectively weighs more than 195 lbs.
2. If an airplane is full of 213 men (and no women or children), each with a 20 lb bag, what is the probability that the total weight is greater than 41535 lbs (the weight limit for the airplane)?
1. With the bag the mean weight is $\mu = 192$. The standard deviation remains the same. $z_{195} \doteq 0.1034$. So $P(z \gt 0.1034) \doteq 0.4588$
2. If the total weight is 41535 lbs, the average weight of the 213 men is 195 lbs. Central limit theorem applies. $\mu = 192$, $\displaystyle{\sigma = \frac{29}{\sqrt{213}} \doteq 1.9870}$. Thus $z_{195} = 2.1282$. So the probability of exceeding the weight limit is $P(z \gt 2.1282) \doteq 0.0167$. | 2021-02-25T08:26:27 | {
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https://math.stackexchange.com/questions/3379668/would-you-use-a-proportion-to-solve-this-problem/3379698 | # Would you use a proportion to solve this problem?
A math class is working with similar figures. Mr. Kole told the class that the surface area of the cones are 180 and 320 square units, and the volume of the smaller cone is 151 cubic units. He challenged his class to find the volume of the bigger cone. What is the volume of the bigger cone?
I tried to use the proportion 320/180 = x/151 to solve the problem. But my answer did not match the correct answer, which is 358 cubic units.
How would you solve this problem?
• Surface areas of similar figures are in proportion to the squares of corresponding lines; volumes are in proportion to the cubes. – saulspatz Oct 3 '19 at 19:27
• Are cones open bases or closed? – Mohammad Riazi-Kermani Oct 3 '19 at 19:27
Since the cones are similar the ratio of their areas is equal to the square of corresponding lines. (For example, the height or the radius of the base.) The ration of those lines is $$a=\sqrt{{320\over180}}=\frac43$$ The ratio of the volume is $$a^3$$ so the volume of the large cone is $$\left(\frac43\right)^3\cdot151$$
Well you could use the formulas for the are and volume of cones to figure everything out and I guess that will be enough to solve the problem. However the problem is much more simple than that.
The area scales as a square of the radius $$A\propto r^2$$ and the volume as the cube $$V\propto r^3$$. In particular $$A^3\propto r^6,\, V^2\propto r^6$$, then the quantity $$\frac{A^3}{V^2}$$ is an invariant for similar cones (or any other figure for that mater). And we have $$\frac{320^3}{x^2} = \frac{180^3}{151^2}$$
This should be alright, however I am getting $$x=358$$ which doesn't seem to be the answer you are getting. Could you please check the numbers?
• Yes. Your answer 358 is right. I made a typo. – orangebull Oct 4 '19 at 3:02
• But why did you put $\frac{A^3}{V^2}$? Shouldn't it have been $\frac{A^2}{V^3}$? – orangebull Oct 4 '19 at 3:04
• Ah sorry, I went a little fast there. $A^3\propto r^6$ and $V^2\propto r^6$ then the ratio of this two quantities doesn't depend on the radius. $\frac{A^3}{V^2}\sim \frac{r^6}{r^6}\sim \text{constant}$. Look that @saulspatz answer is exactly the same argument. – David Jaramillo Oct 4 '19 at 8:23
No, surface area is not directly proportional to volume.
I would try to construct a cone with surface area $$180$$ and volume $$151$$ using the formulas for surface area and volume, in particular keeping track of the ratio between the cone's height and radius.
Then I would construct a cone with surface area $$320$$ and the same ratio between its height and radius. Finally, apply the volume formula to this larger cone.
• All you need is the ratio of the areas. The shape of the cone doesn't matter. – saulspatz Oct 3 '19 at 19:30
• @saulspatz Really? Aren't there infinitely many cones with surface area $320$ but different volumes, depending on their heights and radii? And only one of these will be similar to the smaller cone? – 79037662 Oct 3 '19 at 19:34
• I guess I assumed the cones being similar was relevant information. – 79037662 Oct 3 '19 at 19:35
• It is relevant. That why all you need is the ratio of the areas. – saulspatz Oct 3 '19 at 19:36
• @saulspatz At the first glance I think the ratio alone is not sufficient. – callculus Oct 3 '19 at 19:36
Volumes are proportional to cube and areas to square of linear dimension in similitude.
$$V = k_1 L^3,\, A = k_2 L^2,\, V = k_3 A^{\frac32}$$
Accordingly the larger Volume is
$$V_2 = V_1 \big(\frac{320}{180}\big)^ {1.5} =151\times 64/ 27 \approx 357.926$$
The similarity of two cones imply: $$\frac{R_1}{R_2}=\frac{l_1}{l_2}=\frac{h_1}{h_2} \Rightarrow \frac{R_1}{R_2}=\frac{R_1+l_1}{R_2+l_2}$$ The areas of cones: $$S_1=\pi R_1l_1+\pi R_1^2=180;\\ S_2=\pi R_2l_2+\pi R_2^2=320;\\ \frac{S_1}{S_2}=\frac{\pi R_1(l_1+R_1)}{\pi R_2(l_2+R_2)}=\left(\frac{R_1}{R_2}\right)^2=\frac{180}{320} \Rightarrow \left(\frac{R_1}{R_2}\right)^3=\left(\frac{180}{320}\right)^{3/2}$$ The volumes of cones: $$V_1=\pi R_1^2h_1=151;\\ V_2=\pi R_2^2h_2;\\ \frac{V_1}{V_2}=\frac{\pi R_1^2h_1}{\pi R_2^2h_2}=\left(\frac{R_1}{R_2}\right)^3=\left(\frac{180}{320}\right)^{3/2}=\frac{151}{V_2} \Rightarrow V_2\approx 357.93.$$ | 2020-01-17T19:14:15 | {
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http://math.stackexchange.com/questions/302257/is-ex-the-only-isomorphism-between-the-groups-mathbbr-and-mathbb/302267 | # Is $e^x$ the only isomorphism between the groups $(\mathbb{R},+)$ and $(\mathbb{R}_{> 0},*)$?
If so, how might I be able to prove it?
EDIT: OK, thanks to many answers especially spin's and Micah's explanations. All of the answers were extremely helpful in understanding -- I have accepted Micah's because it seems the most complete, but all answers provide helpful additions/perspectives! I have tried to summarize:
$\phi$ is an isomorphism between the groups if and only if $\phi(x) = e^{f(x)}$ where $f$ is an isomorphism from $(\mathbb{R},+)$ back to $(\mathbb{R},+)$.
Of course there are lots of such $f$, especially when we take the Axiom of Choice.
However, it seems from the answers and Micah's link (Cauchy functional equation) that the only "nice" solutions are $f(x) = cx$ for a constant $c$. It seems that all others must be "highly pathological" (in fact $\{(x,f(x))\}$ must be dense in $\mathbb{R}^2$).
A remaining question is, how strong is the statement
All such isomorphisms have the form $e^{cx}$ for some $c \in \mathbb{R}$
or its negation? Or what is required for each to hold?
One answer seems to be that supposing the reals have the Baire property is sufficient to rule out other solutions (as is assuming every subset of the reals is measurable, assuming the Axiom of Determinacy, and it holds in Solovay's model). For more, see this question, this question, and this mathoverflow question.
-
well you could aslo have $e^{nx}$ for positive $n$. Or $a^nx$. – Joe Tait Feb 13 '13 at 16:35
Ah, good point. I think the question is still interesting when we exclude that case though! – usul Feb 13 '13 at 16:37
As I just added, it also works for any number positive number instead of $e$. Still interesting excluding that though :( – Joe Tait Feb 13 '13 at 16:39
See @GEdgar's comment below. – copper.hat Feb 13 '13 at 16:41
I think the negation of the last one is only a consequence, because it is a statement about the subsets of $\mathbb R$. It cannot imply statements about arbitrary sets.... It follows from the fact that $\mathbb R$ has a $\mathbb Q$-basis. – N. S. Feb 13 '13 at 21:06
show 1 more comment
Let $f:\mathbb{R} \to \mathbb{R}^+$ be an isomorphism. Then $g=\log f$ is an automorphism of $(\mathbb{R}, +)$: that is, it satisfies $$g(x+y)=g(x)+g(y) \, .$$ This is Cauchy's functional equation. The only continuous (or even measurable) solutions are the trivial ones (which correspond to $f(x)=e^{cx}$), but there are also exotic solutions that require some version of the axiom of choice to construct — which would yield similarly exotic $f$s.
-
Very helpful, thanks! My question now (last question?) is how "strong" is the axiom that such solutions exist, or the axiom that they don't? – usul Feb 13 '13 at 21:08
@usul: I'm sorry, I'm not enough of a set theorist to give a good answer to that question. (Though many people here are; you'd probably get good answers if you asked it as a separate question!) – Micah Feb 13 '13 at 21:15
no problem, a search turned up several relevant threads and I have linked to them in my edited question! – usul Feb 13 '13 at 21:20
Once you have an isomorphism, there are plenty. Take two Hamel bases of $\mathbf{R}$, say $(b_i)$ and $(c_i)$, and define $f$ by $$f(b_i) = e^{c_i},$$ extending it by $\mathbf{Q}$-linearity to all of $\mathbf{R}$.
PS Of course I am implicitly using here the argument of @Micah and @spin.
-
And this is where the axiom of choice comes in. – Asaf Karagila Feb 13 '13 at 16:50
@AsafKaragila, precisely! – Andreas Caranti Feb 13 '13 at 16:53
Thanks for the response! – usul Feb 13 '13 at 21:24
$f:\mathbb{R}\rightarrow\mathbb{R}^{+} \text{ satisfying } f(x+y)=f(x)f(y)$, $f$ continuous, then $f(x)=e^{cx}$
Try to solve this problem Any Lie Group Homomorphism from $\mathbb{R}\rightarrow S^1$ is of the form $e^{iax}$ for some $a\in\mathbb{R}$ and every such homomorphism is smooth.
-
(+1). Any ideas if $f$ is discontinous? – Amr Feb 13 '13 at 16:38
Does it really have to be $e$? Why not $a>0$. – Joe Tait Feb 13 '13 at 16:40
Any Lie group Homomorphism is smooth – Bunuelian Trick Feb 13 '13 at 16:43
@Joe Tait $\forall a>0 \exists c\in\mathbb{R} [a=e^c]$ – Amr Feb 13 '13 at 16:46
Thanks for the answer! – usul Feb 13 '13 at 21:26
Suppose that $G \cong H$ as groups and that $f: G \rightarrow H$ is an isomorphism. Then every isomorphism $G \rightarrow H$ is of the form $f \circ g$, where $g: G \rightarrow G$ is an isomorphism. So to determine every isomorphism $G \rightarrow H$, it suffices to find just one and then the rest are given by composing with automorphisms of $G$.
Because $x \mapsto e^x$ is an isomorphism $(\mathbb{R}, +) \rightarrow (\mathbb{R}_{>0}, \cdot)$, every isomorphism $(\mathbb{R}, +) \rightarrow (\mathbb{R}_{>0}, \cdot)$ is of the form $x \mapsto e^{\varphi(x)}$, where $\varphi$ is an automorphism of $(\mathbb{R}, +)$.
-
Dear spin, many thanks for your nice answer! – shankfei Feb 13 '13 at 20:10
Thanks, this was very helpful/clear! – usul Feb 13 '13 at 20:57
$e^{2x}$ is another one. In general if $\phi:\mathbb{R}\rightarrow\mathbb{R}$ is a group isomorphism, then $e^{\phi(x)}$ is another one.
One should also note that if $\phi$ is a group isomorphism between $(\mathbb{R},+)$ and $(\mathbb{R}_{>0},*)$, then $\log(\phi)$ is an automorphism of $(\mathbb{R},+)$. Now the question is to determine the group automorphisms of $\mathbb{R}$ as a group under addition
@Amr: Every isomorphism $\mathbb{R} \rightarrow \mathbb{R}_{>0}$ is of the form $e^{\phi(x)}$, see my answer – spin Feb 13 '13 at 16:53 | 2014-07-24T15:55:56 | {
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# There are 8 books in a shelf that consist of 2 paperback
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There are 8 books in a shelf that consist of 2 paperback [#permalink]
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Updated on: 02 Feb 2014, 22:48
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Question Stats:
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There are 8 books in a shelf that consist of 2 paperback books and 6 hardback books. If 4 books are selected at random without replacement, how many different combinations are there that at least one paperback book is selected?
A. 40
B. 45
C. 50
D. 55
E. 60
Originally posted by pretttyune on 12 Nov 2007, 01:39.
Last edited by Bunuel on 02 Feb 2014, 22:48, edited 3 times in total.
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There are 8 books in a shelf that consist of 2 paperback [#permalink]
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08 Feb 2012, 02:20
19
17
pretttyune wrote:
There are 8 books in a shelf that consist of 2 paperback books and 6 hardback books. If 4 books are selected at random without replacement, how many different combinations are there that at least one paperback book is selected?
A. 40
B. 45
C. 50
D. 55
E. 60
It's almost always better to solve "at least one" combinations/probability questions with {at least one}={total}-{none}.
Total # of selections of 4 books out of 8: $$C^4_8=70$$;
Total # of selections of 4 hardback books out of 6 (so none paperback): $$C^4_6=15$$;
{at least one} = {total} - {none} = 70 - 15 = 55.
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12 Nov 2007, 19:27
8
5
pmenon wrote:
Is the answer above right ?
I was taking the approach of first finding out the number of ways of selecting 4 books out of 8, then finding the number of possibilities where no paperbacks were chosen, and subtracting the two.
If anyone can help with this approach to put down what the actual terms should be, thatd be very helpful to me.
i would go about it this way:
total - all hardback
total = 8C4
hardback = 6C4
8C4 - 6C4
70-15 = 55
##### General Discussion
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12 Nov 2007, 02:12
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5
No of combinations = p(1) + p(2)
= (2C1 * 6C3) + (2C2 * 6C2)
= (2*6*5*4)/(3*2*1) + (6*5)/2
= 40 + 15
= 55
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28 Sep 2009, 03:30
6
4
There are 8 books in a shelf that consist of 2 paperback books and 6 hardback books. If 4 books are selected at random without replacement, how many different combinations are there that at least one paperback book is selected ?
Soln:
No. of combinations which have atleast one paperback is
= Total no. of combinations - Combinations in which there is no paperback
= 8C4 - 6C4
= 55 ways
Other way of doing this is finding out directly
No. of combinations which have atleast one paperback is
= Total no. of combinations that have 2 paperback + total no. of combinations that have 1 paperback
= 2C2 * 6C2 + 2C1 * 6C3
= 15 + 40
= 55 ways
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01 Oct 2009, 23:54
Solution:
Total possible combinations of books = C(8,4)
Total possible combinations without paperback books C(6,4)
Total possible combinations with at least 1 paperback book: C(8,4)-C(6,4)=55
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21 Mar 2011, 04:07
Another approach:
Combinations with the first book:
6C3 = 20
Combinations with the second book:
6C3 = 20
Combinations with both books:
6C2 = 15
Total = 20+20+15 = 55
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### Show Tags
22 Mar 2011, 06:00
PHHH + PPHH
= 2C1 * 6C3 + 2C2 * 6C2
= 2 * 6!/3!3! + 1 * 6!/4!2!
= 2* 6*5*4/3! + 6 * 5/2
= 40 + 15 = 55
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Re: There are 8 books in a shelf that consist of 2 paperback [#permalink]
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08 Feb 2012, 00:59
I came up with a different (more complicated?) approach.
FIRST STEP:
How many possible selections of 4 books from this self including NO paperback?
6/8 * 5/7 * 4/6 * 3/5 = 3/14
SECOND STEP:
Subtract 3/14 from 1 (or 14/14) => 14/14 - 3/14 = 11/14
THIRD STEP:
The only answer choice with 11 as a factor is 55.
Ans. C
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Re: There are 8 books in a shelf that consist of 2 paperback [#permalink]
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08 Feb 2012, 04:30
Do i miss the answer choices?
total combinations: $$C^4_8 = 70$$
No Paperback: $$C^4_6 = 15$$
Total - No Paperback = At least 1 Paperback
70 - 15 = 55 combinations
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Re: There are 8 books in a shelf that consist of 2 paperback [#permalink]
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08 Feb 2012, 09:00
# of selections of books with no condition = 8C4 = 70
# of selections of books with no paperback book = 6C4 = 15
# of selections of books with at least one paperback book = 70 -15 = 55
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Re: There are 8 books in a shelf that consist of 2 paperback [#permalink]
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20 May 2012, 03:36
3
The selection should at least contain one PAPERBACK book. There are a total of 2 PB and 6 HB books are available.
The combinations that at least one PB book will come out are: PHHH & PPHH
1. PHHH = 2C1 * 6C3 = 2 * (5*4) = 40
2. PPHH = 2C2 * 6C2 = 1 * ((6*5) / 2) = 15
In total 40+15 = 55 Ways
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Re: There are 8 books in a shelf that consist of 2 paperback [#permalink]
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25 Jan 2013, 07:20
Hmm....so when I did this I made an error that might have carried over from using combinatorics from probability. I found the total as 8C4 and to find the amount to subtract from it I had (6C4)(2C0). Why don't we multiply by the books not chosen and what's would it mean if we did?
I'm trying to understand a fundamental flaw I made here
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Re: There are 8 books in a shelf that consist of 2 paperback [#permalink]
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25 Jan 2013, 07:33
1
manimgoindowndown wrote:
Hmm....so when I did this I made an error that might have carried over from using combinatorics from probability. I found the total as 8C4 and to find the amount to subtract from it I had (6C4)(2C0). Why don't we multiply by the books not chosen and what's would it mean if we did?
I'm trying to understand a fundamental flaw I made here
You'd get the same result because 2C0=2!/(2!*0!)=2/(2*1)=1: there is one way to choose 0 books from 2 books.
Hope it's clear.
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Re: There are 8 books in a shelf that consist of 2 paperback [#permalink]
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25 Jan 2013, 09:57
# of ways of choosing 4 books from 8 total books = 8C4 = 70
# of ways of only choosing hardback books = 6C4 = 15
Therefore, # of ways of choosing 4 books with at least 1 paperback book = 8C4 - 6C4 = 55 (D)
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Re: There are 8 books in a shelf that consist of 2 paperback [#permalink]
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07 May 2015, 05:03
When I first read the question I understood that we need to find the number of ways to select 4 books with at least 1 of them being a paperback, instead of calculating the number of possible selections, were PHHH = HPHH
Would the answer to this other question be correct, as this might show up on the exam?
2 Paperbacks: 2*1*6*5 = 60
1 Paperback : 1*6*5*4 = 120
= 180
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Re: There are 8 books in a shelf that consist of 2 paperback [#permalink]
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28 Jul 2015, 15:22
First thing to notice when you look at this question: we're dealing with permuations/combinations.
Which to use? Ranking/order doesn't matter here - which means we use combinations.
Next thing to notice: at least one - i don't know anyone who likes these questions!
Strategy should be TOTAL - NONE = AT LEAST 1 (our desired solution)
TOTAL = 8c4 (8 books 4 being selected) = 70
NONE = 6c4 (Only 6 books because we aren't including paperback, 4 being selected) = 15
70-15 = 55
CHOOSE D
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Re: There are 8 books in a shelf that consist of 2 paperback [#permalink]
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07 Mar 2016, 03:52
a method to solve this problem
so there are 2 possibilities
either : 1 PB and 3 HB = 2*6*5*4/3! = 40
OR : 2 PB and 2 HB = 2*1*6*5/2!*2! = 15
so we have to add both the possibilities to get the final answer
40+15=55
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There are 8 books in a shelf that consist of 2 paperback [#permalink]
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28 Jul 2016, 05:05
Bunuel wrote:
pretttyune wrote:
There are 8 books in a shelf that consist of 2 paperback books and 6 hardback books. If 4 books are selected at random without replacement, how many different combinations are there that at least one paperback book is selected?
A. 40
B. 45
C. 50
D. 55
E. 60
It's almost always better to solve "at least one" combinations/probability questions with {at least one}={total}-{none}.
Total # of selections of 4 books out of 8: $$C^4_8=70$$;
Total # of selections of 4 hardback books out of 6 (so none paperback): $$C^4_6=15$$;
{at least one}={total}-{none}=70-15=55.
Hi Bunuel,
Using PB: Paperbacks ; HB : Hardbacks
If the question were - "How many possible selections of 4 books include both the PB" then is why is my following solution wrong ?
Total possibility .: C(8,4)
Possibility of 2 PB = Total possibility - possibility of no PB => C(8,4)- C(6,4) [C(6,4) -> removing 2 PB from the lot will leave only 6 books. So selecting 4 books from 6].
So answer is C(8,4)- C(6,4) = 70-15 = 55
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Re: There are 8 books in a shelf that consist of 2 paperback [#permalink]
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27 Aug 2016, 02:23
1. Method :- 8c4 - 6c4 = (Total - All four book selected from 6 hard copy) - This will give atleast one paper back selected.
2.Method :- 6c2*2c2 + 6c3*2c1 = (PPHH) + (PHHH)
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Sumit kumar
If you like my post Kudos Please
Re: There are 8 books in a shelf that consist of 2 paperback &nbs [#permalink] 27 Aug 2016, 02:23
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Display posts from previous: Sort by | 2019-01-22T14:51:05 | {
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http://math.stackexchange.com/questions/179875/what-is-a-metric-for-mathbb-q-in-the-lower-limit-topology | What is a metric for $\mathbb Q$ in the lower limit topology?
A useful source of counterexamples in topology is $\mathbb R_\ell$, the set $\mathbb R$ together with the lower limit topology generated by half-open intervals of the form $[a,b)$. For example this space is separable, Lindelöf and first countable, but not second countable. So in particular, this implies that $\mathbb R_\ell$ is not metrizable.
Now Urysohn's metrization theorem says that any regular, second countable space is metrizable. Applying this to $\mathbb Q_\ell\subset \mathbb R_\ell$ gives the (for me rather counterintuitive) fact that there must exist a metric on $\mathbb Q$ inducing the lower limit topology. Which leads to the question:
What is a concrete example of a metric inducing the lower limit topology on $\mathbb Q_\ell$?
It certainly should be possible to go through the steps of the Urysohn metrization theorem to give a (semi-)concrete embedding of $\mathbb Q_\ell$ in $\mathbb R^\omega$ (for instance picking an enumeration of the rationals and constructing functions which separate points from closed sets, then look at the cartesian product of these maps) and then give the metric on $\mathbb Q_\ell$ in terms of this embedding.
But this is really not what I'm looking for, here. I'd be much more interested in a simple metric for $\mathbb Q_\ell$.
Because "simple" is not well-defined, I will definitely also welcome answers which exhibit a metric, but which I would not consider to be simple. On the other hand, if you see a reason for why there may just be no "simple" metrics, please feel free to point this out as well.
Thanks!
-
I thought $\mathbb{Q}_{\ell}$ was refering to the $\ell$-adic numbers. – M Turgeon Aug 7 '12 at 13:11
Why is $\mathbb{Q}_l$ second countable? – Rob Arthan Aug 7 '12 at 19:54
@RobArthan: Because the open sets of the form $[p,q)$ for $p,q\in \mathbb Q$ form a basis of the topology. – Sam Aug 7 '12 at 20:06
@MTurgeon: It seems you are not alone in being mislead by the old title. Thanks for pointing it out! My intention was not to mislead anyone, of course, so I have changed it now. – Sam Aug 7 '12 at 20:09
@Sam L. Thank you! I was being stupid. – Rob Arthan Aug 7 '12 at 21:04
Let $\nu:\mathbb{Q} \to \mathbb{N}$ be an enumeration of $\mathbb{Q}$. Then $\displaystyle d(x,y):=\sum_{\min(x,y) < r \le \max(x,y)} 2^{-\nu(r)}$ will do.
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@tomasz: Not so. Notice, for instance, that for any $x < 0$, we have $d(x,0) \ge 2^{-\nu(0)}$. – Nate Eldredge Aug 7 '12 at 21:26
@NateEldredge: good point. :) This is a nice example. – tomasz Aug 7 '12 at 21:31
A remark: If we let $f(x) = \sum_{r \le x} 2^{-\nu(r)}$, then $d(x,y) = |f(x) - f(y)|$. In particular, $\mathbb{Q}$ with this metric is isometric (and in particular homeomorphic) to the set $\{f(q) : q \in \mathbb{Q}\} \subset [0,1]$ with the Euclidean metric. – Nate Eldredge Aug 7 '12 at 21:47
This is very nice. Thank you! So given an interval $[p,q)$, the $\epsilon$-ball $B_\epsilon(p)$ will be contained in $[p,q)$ if $\epsilon = \min\{2^{-\nu(p)},2^{-\nu(q)}\}$. Brilliant. – Sam Aug 8 '12 at 5:11
EDIT: This one should work out.
Here's another option: write your rationals as "mixed fractions," that is as their integer floor plus a fractional part and define $$d\left(a\frac{p}{q},b\frac{r}{s}\right)=\begin{cases}|a-b|,& a\neq b\\ d'\left(\frac{p}{q},\frac{r}{s}\right), &\textrm{otherwise} \end{cases}.$$
Define the distance between pure fractions separately. Assume WLOG that $\frac{p}{q}\leq \frac{r}{s}$. Then set $$d'\left(\frac{p}{q},\frac{r}{s}\right)=\max\left(\left|\frac{p}{q}-\frac{r}{s}\right|,\frac{1}{m}\right), \frac{k}{m} \in \left(\frac{p}{q},\frac{r}{s}\right]$$ We see, for instance, that this gives $\left[\frac{1}{2},\frac{5}{6}\right)$ open, since the distance from anything smaller than $\frac{1}{2}$ to $\frac{1}{2}$ is $\frac{1}{2},$ but each distance from $\frac{1}{2}$ to something less than $\frac{5}{6}$ is no more than $\frac{1}{3}$.
The symmetry and homogeneity axioms are immediate. Let's consider the triangle inequality. The only case in which it doesn't follow from that for the standard absolute metric is when we must show $d'\left(\frac{p}{q},\frac{r}{s}\right)+d'\left(\frac{r}{s},\frac{t}{u}\right)\geq d'\left(\frac{p}{q},\frac{t}{u}\right)$ and $d'\left(\frac{p}{q},\frac{t}{u}\right)=\frac{1}{m}$. But if the inequality failed, we'd have to have the interval $\left(\frac{p}{q},\frac{t}{u}\right]$ contain something with denominator $m$ bigger than any denominator in the intervals on the left-hand side-which is absurd, since depending on the ordering we have one of $\left(\frac{p}{q},\frac{t}{u}\right]=\left(\frac{p}{q},\frac{r}{s}\right]\cup\left(\frac{r}{s},\frac{t}{u}\right], \left(\frac{p}{q},\frac{t}{u}\right]\subset \left(\frac{p}{q},\frac{r}{s}\right],$ or $\left(\frac{p}{q},\frac{t}{u}\right]\subset\left(\frac{r}{s},\frac{t}{u}\right]$.
So we have a metric. To get some half-open interval $[x,y)=\left[a\frac{p}{q},b\frac{r}{s}\right)$, take the union of all the $[m,m+1) \subset [a,b)$. Then for the least such $m$, construct $[x,m)$ by the infinite union $\bigcup_{i=0}^\infty\left[a\frac{p}{q+i},a\frac{p}{q+i}+\frac{1}{q+i+1}\right)$, getting all these intervals by the argument above about $\left[\frac{1}{2},\frac{2}{3}\right)$. Get $\left[b,b\frac{r}{s}\right)$ as the ball of radius $\frac{r}{s}$ around $b$, union the three pieces together, and we're done.
-
Do you mean $a\neq b$ in the first formula? – Grumpy Parsnip Aug 7 '12 at 21:10
Yep, thanks @JimConant. – Kevin Carlson Aug 7 '12 at 21:11
@tomasz: You're right, this won't do it. If it would work for the rationals, it would work for the reals. – Kevin Carlson Aug 7 '12 at 21:18
This seems to work. It's nowhere near as elegant as Alexander Shamov's answer, though, especially considering the many cases (some of which you omitted) with triangle inequality... Still, kudos on your determination. :) – tomasz Aug 7 '12 at 22:40
@tomasz, I agree-Alexander's solution just wasn't intuitive to me, so I thought it wouldn't be completely useless to share one with a few more of the strings still showing. – Kevin Carlson Aug 7 '12 at 23:05 | 2015-05-29T00:24:40 | {
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http://mathhelpforum.com/advanced-math-topics/74420-incidence-geometry-print.html | # Incidence geometry
• February 18th 2009, 09:13 PM
namelessguy
Incidence geometry
This seems like hard problem. Anyone has any suggestion how to tackle it?
Show that in a finite incidence geometry, the number of lines is greater than or equal to the number of points.
• February 18th 2009, 09:40 PM
ThePerfectHacker
Quote:
Originally Posted by namelessguy
This seems like hard problem. Anyone has any suggestion how to tackle it?
Show that in a finite incidence geometry, the number of lines is greater than or equal to the number of points.
Between any two distinct points there exists exactly one line. This is one of the axioms of incidence geometry. Now if there are $n$ points (and say $n\geq 2$) then the number of pairs that we can form is equal to ${n\choose 2}$ while the number of lines is $n$.
EDIT: Mistake
• February 18th 2009, 09:58 PM
namelessguy
Quote:
Originally Posted by ThePerfectHacker
Between any two distinct points there exists exactly one line. This is one of the axioms of incidence geometry. Now if there are $n$ points (and say $n\geq 2$) then the number of pairs that we can form is equal to ${n\choose 2}$ while the number of lines is $n$.
Did you mean that the number pairs we can form is equal to ${n\choose 2}$ which is the number lines, while the number of points is only $n$? I did learn the axiom, and thanks a lot for your help TPH. It seems much easier.
• February 19th 2009, 10:23 PM
namelessguy
@TPH, actually, I tried to draw some examples for some n points, but it appears that the number lines can be less that ${n\choose2}$. I guess it could be that ${n\choose2}$ is the maximum number of lines that we can get from n points. If this is the case, then I'm stuck again(Speechless)
• February 22nd 2009, 12:12 PM
Plato
Quote:
Originally Posted by namelessguy
I tried to draw some examples for some n points, but it appears that the number lines can be less that ${n\choose2}$. I guess it could be that ${n\choose2}$ is the maximum number of lines that we can get from n points.
That is truly the case.
Example: The Pappus finite geometry. In this geometry there are nine points and nine lines with 3 points on each line and 3 lines on each point.
Thus this problem depends on the set of axioms in use.
• March 13th 2009, 06:40 PM
George Jennings
Incidence geometry: #lines >= #points
In an incidence geometry, two points are connected by exactly one line, there are (at least) three noncollinear points, and every line contains at least two points. In this context the de Bruijn-Erdos theorem (1948) says that there are at least as many lines as points. There's a proof in the text Combinatorics of Finite Geometries by Lynn Margaret Batten. One can also find proofs on the net by searching for "de Bruijn-Erdos theorem".
George | 2014-03-14T19:50:31 | {
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https://math.stackexchange.com/questions/697818/factoring-x3-x2-x-1 | Factoring $x^3 - x^2 - x + 1$
I'm familiar with basic forms of polynomial factoring like completing the square or factoring (e.g. finding that $x^2+x-6=(x+3)(x-2)$, but I'm currently working on integration of rational functions by partial fractions and I find myself encountering the factorizations of polynomials unlike those I've worked with before.
Given: $x^3 - x^2 - x + 1$
When I first look at this, I'm unsure of where to start. It's very unlike the example above. It's degree $3$, for one thing, and it has $4$ terms.
When I look ahead to the solution, the book says $f(1) = 0$, so we know that $(x-1)$ is a factor. Okay, now that they mention that, I can see it fits, but what?! How did they come up with that? Is that a general rule? Does anyone have a link to where I can learn about it?
Further, once it's known that $(x-1)$ is a factor, how do I then factor an $(x-1)$ out of a longer polynomial like this?
In other words, how does $x^3 - x^2 - x + 1 = (x-1)(x^2-1)$?
• polynomial division. – Alexander Grothendieck Mar 3 '14 at 16:59
• $f(x) = x^{3} - x^{2} - x + 1$ so by substituting 1 in for all of the $x$ values you obtain $f(1) = 0$ so you have $x - 1$ is a factor. Then use this example purplemath.com/modules/polydiv2.htm for the other factor – Ahrz Mar 3 '14 at 17:00
• – Alraxite Mar 3 '14 at 17:02
Answer to your 1st question: if a polynomial is degree $n$, it can have up to $n+1$ terms, so having $4$ terms in a degree $3$ polynomial is nothing out of the ordinary.
Answer to your 2nd question: finding a factor of a polynomial of degree $3$ or more is usually just guesswork. In fact, it's been proven that it's IMPOSSIBLE to find a formula to find all solutions of a polynomial with degree of $5$ or above. And unless you're willing to memorize this gruesome formula for polynomials of degree 3 (and let's not even talk about degree 4...), guesswork is your best friend.
• Why memorize, if computer algebra packages will factor for you (if there are nice factors)? – vonbrand Mar 3 '14 at 17:56
\begin{align}x^3-x^2-x+1&=x^3-x^2-(x-1)\\&=x^2(x-1)-(x-1) \\&=(x-1)(x^2-1)\\&=(x-1)[(x-1)(x+1)]\\&=(x+1)(x-1)^2\end{align}
• I have edited your post for readability. In the future, to make multiple lines in MathJax, put \\ after the end of one line and &= in place of = to align the equations like above. – Cameron Williams Mar 3 '14 at 17:45
Three answers and nobody has mentioned the Factor Theorem!
Factoring cubics can be tricky. But, if you know a few tricks, then it might not be so bad.
One trick for this cubic: there appears to be a pattern in the coefficients: $+1,-1,-1,+1$, and this sums to zero, so it's logical to look at $\pm 1$ as a root.
More generally, you can use the rational root theorem to impose conditions on any rational root of the polynomial. In particular, the rational root theorem says that if $\frac{p}{q}$ is a root of the polynomial, then $q$ must divide the leading coefficient ($1$), and $p$ must divide the unit coefficient (also $1$). Therefore, if the polynomial has a rational root, $1$ is a candidate ($-1$ might also work).
Once you impose conditions on nicely-behaved roots, simply try to factor it out by polynomial division!
• In fact, $1$ is a root of any polynomial whose coefficients sum to zero! – Emily Mar 3 '14 at 17:04
One thing that can be done for polynomials of degree 3 is that you can estimate the roots by differentiating the polynomial and equationg that expression to 0. That will hopefully help you to find a domain in which one root lies. For example here, $$\frac d{dx}(x^3-x^2-x+1)=3x^2-2x-1=0$$$$=>(3x+1)(x-1)=0 \\ =>x=1 \space and \space x=\frac {-1}3.$$ Therefore $\frac {-1}3 \le \alpha\le 1$. Then you may try integer values in the domain that satisfy the polynomial. In this case it is 1 and since 1 is also a root of the derivative, hence it is a double root. Now you can easily find the other root that is -1.
There's some theorem (rational root theorem/ factor theorem?) that says that for all rational factors of the form p/q, p divides into the constant in your polynomial and q divides into the leading coefficient. That leaves to with 1 and -1. Testing 1, first, we find it's a root. -1 is also a root. Using polynomial division it becomes easier to test these numbers. Since it's a cubic it has three roots - it turns out that 1 is a double root. | 2019-11-15T20:40:06 | {
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https://math.stackexchange.com/questions/3023292/right-triangle-geometry-problem | Right triangle geometry problem
Right triangle $$\Delta ABC$$ ($$\angle ACB=90°$$). The following is constructed: from point $$C$$ altitude $$CD$$, angle bisector $$CL$$ of $$\angle ACB$$, angle bisector $$DK$$ of $$\angle ADC$$, angle bisector $$DN$$ of $$\angle BDC$$.
$$D, L$$ lie on $$AB$$, $$K$$ lies on $$AC$$, $$N$$ lies on $$BC$$.
Prove that $$C, K, L, D, N$$ lie on same circle and prove that $$|KN|=|CL|$$
I think that I need to do something with quadrilateral $$CKDN$$. I got that $$|KN|=\sqrt{|CK|^2+|CN|^2}=\sqrt{|DK|^2+|DN|^2}$$ $$\angle CKD+\angle CND=180°$$ I also tried to express sides with the angle bisector theorem, but I don't know how to continue / what I need to solve this. How can I solve this problem?
• No they aren't. – Pero Dec 3 '18 at 12:54
First, since $$CD \perp AB$$ and $$DK$$ and $$DN$$ are angle bisectors to the right angles $$\angle \, ADC$$ and $$\angle \, BDC$$, then $$\angle \, KDN = \angle \, KDC + \angle \, NDC = 45^{\circ} + 45^{\circ} = 90^{\circ}$$ However, $$\angle \, KCN = 90^{\circ}$$ so the quadrilateral $$CKDN$$ is inscribed in a circle.
Next, prove that $$KL\, || \, CB$$ and $$NL\, || \, CA$$ using the properties of angle bisectors and the similarity between triangles $$ABC, ACD$$ and $$BCD$$. Indeed, since $$DK$$ is a bisector of the angle at vertex $$D$$ of triangle $$\Delta \, ADC$$, we apply the theorem that $$\frac{AK}{KC} = \frac{AD}{DC}$$ But triangles $$\Delta \, ACD$$ is similar to $$\Delta \, ABC$$ so $$\frac{AD}{DC} = \frac{AC}{CB}$$ so $$\frac{AK}{KC} = \frac{AC}{CB}$$ By the fact that $$CL$$ is an angle bisector of the angle at vertex $$C$$ of triangle $$\Delta\, ABC$$ we have that $$\frac{AC}{CB} = \frac{AL}{LB}$$ so consequently $$\frac{AK}{KC} = \frac{AL}{LB}$$ which by Thales' intercept theorem implies that $$KL \, || \, CB$$. Analogously, one can show that $$NL \, || \, CA$$.
Then quad $$CKLN$$ is a rectangle, so $$KL =NL$$ as diagonals in a rectangle. Therefore the point $$L$$ also lies on the circumcircle of quad $$CKDN$$ and $$KN$$ and $$CL$$ are diameters of the said circle.
• What is $\angle ALLB$? (why not just $\angle ALB$?) – Pero Dec 3 '18 at 23:16
• Oh I see. I think you also have a typo at the beginning (3rd row). I think it should be $CKDN$ is inscribed in circle (not $CKDL$) – Pero Dec 4 '18 at 0:04 | 2019-09-15T18:11:23 | {
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https://math.stackexchange.com/questions/2412828/solve-int-02-sqrtx-sqrtx-sqrtx-dotsb-dx | # Solve $\int_0^2 \sqrt{x+\sqrt{x+\sqrt{x+\dotsb}}}\,dx$
I haven't seen this question, but if someone has, it would be very appreciated if you could send a link!
I've been very interested in the MIT Integration Bee, and one question that stood out to me was:
$$\int_0^2 \sqrt{x+\sqrt{x+\sqrt{x+\dotsb}}}\,dx$$
I tried rewriting a couple of ways to simplify it, but nothing seemed to help. According to the official MIT Integration Bee website, the answer is $$\frac{19}{6}$$
Thanks!
• Can you identify the infinite nesting radical with a closed form? – Sangchul Lee Sep 1 '17 at 2:46
• Hint: If $y=\sqrt{x+y}$, then what is $y$ in terms of $x$? – Shalop Sep 1 '17 at 2:49
• @Shalop Convergence still needs to be justified. – MathematicsStudent1122 Sep 1 '17 at 2:50
• @MathematicsStudent1122 Hence a comment, not an answer. – Shalop Sep 1 '17 at 2:53
To formally justify this integration, we can do the following. Put $$f_1(x) = \sqrt x$$ ,and for each $$n\ge 1$$, put $$f_{n+1}(x) = \sqrt{x + f_{n}(x)}$$. We want to show that $$f(x) = \lim_{n\to\infty}f_n(x)$$ exists and is integrable on $$[0,2]$$. We will show that $$f(x)$$ is continuous on $$(0,2]$$, and has a discontinuity at $$x = 0$$, where $$f(0) = 0$$. This is sufficient because a function $$f\colon[0,2]\to\mathbf R$$ is Riemann integrable if and only if it is bounded and continuous almost everywhere. Finally, by "correcting" the value of $$f$$ at $$x=0$$ to produce a continuous function $$\tilde f$$ on $$[0,2]$$, we can evaluate the integral $$\int_0^2f(x)\,\mathrm dx$$ by applying the fundamental theorem of calculus to integrate $$\int_0^2\tilde f(x)\,\mathrm dx$$.
For each $$x\in (0,2]$$, put $$x_n = f_n(x)$$. We claim that $$x_n\to \frac{1}{2} + \frac{1}{2}\sqrt{4x+1}$$. First we will justify the convergence of the sequence $$(x_n)$$ using the monotone convergence theorem and then use algebraic limit laws to deduce its limit rigorously.
Consider the sequence $$b_{n+1} = \sqrt{2 + b_n}$$, where $$b_1 = \sqrt 2$$. Clearly $$b_1 \leqslant 2$$. Suppose that $$b_n\leqslant 2$$; then $$\sqrt{2 + b_n} \leqslant \sqrt{2 + 2} = 2$$. By induction, $$b_n \leqslant 2$$ for every $$n$$. Thus the sequence $$(b_n)$$ is bounded above by $$2$$. The sequence $$(b_n)$$ is clearly monotone increasing. Hence $$(b_n)$$ converges.
Since $$(x_n)$$ is monotone and bounded above by $$\lim_{n\to\infty}b_n$$, it converges, so set $$L = \lim_{n\to\infty}x_n$$. Since $$(x_n)$$ converges to $$L$$, every subsequence converges to $$L$$, so by continuity of the map $$y\mapsto y^2$$, we have $$L^2 = \big(\lim_{n\to\infty}x_{n+1}\big)^2 = \lim_{n\to\infty}x_{n+1}^2 = \lim_{n\to\infty}x + x_n = x + L.$$ Therefore $$L$$ is a root of the polynomial $$y\mapsto y^2 - y - x$$, and we can conclude that $$L = \frac{1}{2} + \frac{1}{2}\sqrt{4x+1}.$$ (Note that $$L$$ could not be the negative root since $$x_n \geqslant 0$$ for every $$n\geqslant 1$$ and every $$x\in(0,2]$$.)
If $$x = 0$$, then $$x_n = 0$$ for each $$n$$, so all told, $$f(x) = \begin{cases} \frac{1}{2} + \frac{1}{2}\sqrt{4x+1}, &\text{if x\in (0,2],}\\ 0, &\text{if x = 0.} \end{cases}$$
Thus we have proved that the integrand $$f(x)=\sqrt{x+\sqrt{x+\dotsb}}$$ is bounded and continuous almost everywhere on $$[0,2]$$, so it is Riemann integrable on $$[0,2]$$. The function $$\tilde f\colon [0,2]\to\mathbf R$$ defined by $$\tilde f(x) = \frac{1}{2} + \frac{1}{2}\sqrt{4x+1},\quad\text{for each x\in[0,2],}$$ is continuous on $$[0,2]$$ and agrees with $$f$$ almost everywhere, so $$\int_0^2 f(x)\,\mathrm dx = \int_0^2\tilde f(x)\,\mathrm dx$$. Since $$\tilde f(x)$$ is continuous on the entire closed interval $$[0,2]$$, its integral can be evaluated via the fundamental theorem of calculus. Since $$F(x) = \frac{x}{2} + \frac{(4x+1)^{3/2}}{12}$$ satisfies $$F'(x) = \tilde f(x)$$ for each $$x\in [0,2]$$, by the fundamental theorem of calculus, \begin{align*} \int_0^2\tilde f(x)\,\mathrm dx = F(2)-F(0) = \bigg[\frac{2}{2} + \frac{(4\cdot 2+1)^{3/2}}{12}\bigg] - \bigg[\frac{0}{2}+\frac{(4\cdot0+1)^{3/2}}{12}\bigg]=\frac{19}{6}, \end{align*} as desired.
• Just to elaborate, the sequence $b_n$ is monotone increasing because $u \leq \sqrt{u+2}$ whenever $|u| \leq 2$ – Shalop Sep 1 '17 at 17:32
Informally, if we set the nested radical equal to $y$ then an observation to make is that we have $$y = \sqrt{x+y}.$$ This leads us to $$y^2-y-x=0.$$ Looking at the positive solution to this, we have $$y = \frac{1}{2}(\sqrt{4x+1}+1).$$ So, we integrate the $y$ that we just found over the interval from 0 to 2 to achieve $\frac{19}{6}.$
Note that this is very informal. I suspect that justifying convergence in the middle of MIT's integration bee is not necessary.
• For the sake of justification, how would you do it? – John Lou Sep 1 '17 at 3:00
• @JohnLou I don't know for sure but my guess is that one would set up a recursive formula $x_{n+1}=\sqrt{x+x_n}$ and use monotone convergence. Perhaps later, I could try to come up with a full argument. I'm not entirely sure what the initial value would be in this case. – user328442 Sep 1 '17 at 3:03
• How can you justify that we have to take the positive quadratic (as opposed to $1-\sqrt{4x+1}$? – John Lou Sep 1 '17 at 3:06
• I assumed that to be the case since the interval of interest contains value of x from 0 to 2. If we consider these values then y ranges between 0 and -1, which doesn't seem appropriate given a nested square root. – user328442 Sep 1 '17 at 3:08 | 2019-09-18T23:48:09 | {
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https://math.stackexchange.com/questions/1427264/conditional-probability-balls-in-an-urn/1427287 | # Conditional probability - Balls in an urn
I'm quite new to the whole probability branch of math, have been doing some problems, and I've come across a challenging one (in my inexperienced opinion):
There are 50 balls in an urn, 20 red, 30 white. If they are randomly being drawn without being returned, what's:
1. The probability to have a white ball drawn in the 3rd draw
2. The probability to have a red ball drawn in the 7th draw
3. The probability to have drawn a red ball in the 7th draw if it's known that a white ball was drawn in the 3rd draw, and vice versa.
So, I've had 2 problems I just couldn't wrap my head around, and those are:
• How to calculate the probabilities for 1. and 2. I went brute force for the probability of a white ball being drawn on the 3rd draw, made a total probability tree (diagram), summed up all the good outcomes, and got $\frac35$, which would be the same as if it was being drawn in the first draw. Then I checked the same for the 2nd draw and the outcome was the same, so I could assume that the answer to question 2. would be $\frac25$, but I'm still unsure as to why, and would really appreciate an explanation to this.
• Even with those 2 solved, not quite sure how to answer question 3.
I'd be really thankful for all and any kind of explanations to these 2 problems.
• 1. Are any of the other balls+urn questions helpful? 2. What exact condition is meant by "vice versa" in part 3? – shoover Sep 8 '15 at 22:16
This is a good question. Developing an understanding for why the probability is the same as for the first draw will help you find elegant solutions to probability problems that would otherwise require tedious calculations (like the ones you just performed).
To my mind, the most immediate reason that the probability of drawing a white ball must be the same on the third draw as on the first is that, since the balls are by definition identical (except for their colour, which doesn't influence their chances of being drawn), on any given draw all balls must have the same probability of being drawn. But that means that the probabilities of all events related to a single draw (including the event of drawing a white ball) must be the same for all draws.
For your third question, the fact that a white ball was drawn in the third draw simply eliminates one of the white balls, so the probability that a red ball is drawn in any of the other draws (before or after the third draw) is $20/49$. | 2019-10-17T02:42:41 | {
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https://math.stackexchange.com/questions/542634/31-331-3331-33331-333331-3333331-33333331-are-prime | # 31,331,3331, 33331,333331,3333331,33333331 are prime
31,331,3331, 33331,333331,3333331,33333331 are prime. This law can continue it? Will there emerge a composite number? Without using a computer how to judge.
• It won't continue for long. $n \mapsto 10\cdot n + 21$ has good chances to become a multiple of almost any prime dividing neither of $10$ and $21$. Oct 28, 2013 at 10:02
• You might be interested in "suffix primes": johndcook.com/blog/2013/03/12/a-suffix-prime - there are a finite number of them. Oct 28, 2013 at 10:06
• Here's a generalization of this problem on mathoverflow: mathoverflow.net/questions/50071/prepending-strings-to-primes Oct 28, 2013 at 13:18
• Oct 28, 2013 at 18:19
• If you have a pattern or formula you can be near certain it will not generate only primes since there is no simple method like this known to generate primes. So given the choice of trying to prove the pattern continues or look for a counterexample, look for the counterexample. Jul 30, 2014 at 18:11
333333331 is not prime; it is divisible by 17. This does not require a computer. Euler did calculations like this all the time.
What's more, in your sequence 31, 331, 3331, 33331, …, every 15th number is divisible by 31.
Proof: An noted in lab bhattacharjee's answer, the sequence has the form $$a_n = \frac{10^{n+1}-7}{3}$$ Now, 15 is the multiplicative order of $10 \pmod{31}$, so $$a_{15k+1} = \frac{10^{15k+2}-7}{3} \equiv \frac{10^2-7}{3} \equiv 0 \pmod{31}.$$
It has been proven that for all sequences that look like $ab$, $abb$, $abbb$, $abbbb$, … or $ab$, $aab$, $aaab$, $aaaab$, … where the $a$ and $b$ are digits, that periodically the numbers in the sequence are divisible by the first number $ab$.
As an easy exercise, show that in the sequence 11, 111, 1111, 11111, …, that every second term is divisible by 11.
HINT:
$$\underbrace{33\cdots33}_{n \text{ digits}}1=10\frac{10^n-1}3+1=\frac{10^{n+1}-7}3$$
We need to find a prime divisor$(p)$ of $\frac{10^{n+1}-7}3$
Observe that $p>11$
Another way will be to find $p$ such that $10$ is a primitive root
• As $3$ is a primitive root of $17,$ so will be $3^3=27\equiv10\pmod{17}$ and find this dbfin.com/2012/04/… Oct 28, 2013 at 10:16
A little bit off-topic answer:
You can't assume that statements that hold for small number will hold for big number, especially for sequences with primality, because even with computer technology the prime number are one of the biggest mysteries of the maths world. There is one brilliant saying, I can't remember exactly, but one great mathematician (I think it was Polya) once said "An intuition isn't a proof. Which is right for you case.
Also one another good quote is "There'll never be enough small number to prove that something holds generally.". This two quotes describe your problem, even if every number till the $1000^{th}$ one is prime, that doesn't prove that the next one isn't. Luckily in this problem the break happens at the $8^{th}$ number, and as other users proved, there are infinite amount of composite numbers in the sequence.
And at last I think you might be interested in the piece "Law of Small Numbers" by Richard Guy. There are lot of examples like yours, actually this example is included into the book. You can see that the conjectures are proven for some sequences, but the most of them are false. And for some examples the first off-example happens quite late.
There is a note in OEIS sequence A123568 Prime numbers of the form $$\frac{(10^n - 7)}{3}$$. which says:
Note that each $$n$$ from 2 to 8 gives primes, but after that the $$n$$ that correspond to primes are progressively further apart. Singh (1997) gives this as an example of why mathematicians don't trust a preponderance of evidence as proof: in the 17th century, when factoring numbers with as few as eight digits wasn't as easy as it is today, the pattern suggested that all numbers of this form are prime. - Alonso del Arte, Nov 11 2012 | 2022-05-21T09:59:52 | {
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https://boredofstudies.org/threads/induction-question.389810/ | # induction question (1 Viewer)
#### tutor01
##### Member
Can anyone prove by induction that
n(x+a)^(n-1) = sum_{r=0}^{n} r . nCr . x^(r-1) . a^(n-r) for n=1,2,3,...
Thanks.
#### ultra908
##### Member
Uh can you send a photo or write in Latex because your equation is a bit hard to understand.
#### Drdusk
##### π
Moderator
Uh can you send a photo or write in Latex because your equation is a bit hard to understand.
$\bg_white \text{Prove}$
$\bg_white n(x+a)^{n-1} = \sum_{r = 0}^{n}r{n\choose r}x^{r-1}a^{n-r}$
#### tutor01
##### Member
$\bg_white \text{Prove}$
$\bg_white n(x+a)^{n-1} = \sum_{r = 0}^{n}r{n\choose r}x^{r-1}a^{n-r}$
Yes that's the one.
#### harrowed2
##### Member
Here is my working for that question. I have left the n=1 step out to allow the solution to fit on one page.
#### Attachments
• 3 MB Views: 38
#### Drongoski
##### Well-Known Member
Love your very nice handwriting too, harrowed2.
Thank you!
#### CM_Tutor
##### Active Member
Unless required to prove this by induction, this is an excellent example of where an alternative method is much easier.
This result is the direct application of differentiation to the binomial theorem.
The binomial theorem states that $\bg_white (x+a)^n = \sum_{r=0}^n {^n C_r} a^{n - r} x^r = {^n C_0} a^n + {^n C_1} a^{n - 1} x + {^n C_2} a^{n - 2} x^2 + ... {^n C_n} x^n$
Differentiating with respect to x yields: $\bg_white n (x+a)^{n-1} \times 1 = 0 + {^n C_1} a^{n - 1} + {^n C_2} a^{n - 2} 2x + ... {^n C_n} nx^{n-1} = \sum_{r=1}^n {^n C_r} a^{n - r} r x^{r-1}$
Noting the term in the sum is zero when $\bg_white r = 0$, we have: $\bg_white n (x+a)^{n-1} = \sum_{r=0}^n r {^n C_r} x^{r-1} a^{n - r}$ as required
Something to bear in mind here: Questions sometimes use the word "otherwise", as in "Using induction or otherwise, prove ...". The word "otherwise" pretty much always indicates that either there is an alternative method that is much quicker, or that there is an alternative approach that will lead to a huge mess and you should avoid it. A question with a preferable "otherwise" approach offers the chance to pick up time in an exam, which you can then use on other questions. It is worth stopping to consider when the word "otherwise" appears and asking if a better way is available.
#### harrowed2
##### Member
Unless required to prove this by induction, this is an excellent example of where an alternative method is much easier.
This result is the direct application of differentiation to the binomial theorem.
The binomial theorem states that $\bg_white (x+a)^n = \sum_{r=0}^n {^n C_r} a^{n - r} x^r = {^n C_0} a^n + {^n C_1} a^{n - 1} x + {^n C_2} a^{n - 2} x^2 + ... {^n C_n} x^n$
Differentiating with respect to x yields: $\bg_white n (x+a)^{n-1} \times 1 = 0 + {^n C_1} a^{n - 1} + {^n C_2} a^{n - 2} 2x + ... {^n C_n} nx^{n-1} = \sum_{r=1}^n {^n C_r} a^{n - r} r x^{r-1}$
Noting the term in the sum is zero when $\bg_white r = 0$, we have: $\bg_white n (x+a)^{n-1} = \sum_{r=0}^n r {^n C_r} x^{r-1} a^{n - r}$ as required
Something to bear in mind here: Questions sometimes use the word "otherwise", as in "Using induction or otherwise, prove ...". The word "otherwise" pretty much always indicates that either there is an alternative method that is much quicker, or that there is an alternative approach that will lead to a huge mess and you should avoid it. A question with a preferable "otherwise" approach offers the chance to pick up time in an exam, which you can then use on other questions. It is worth stopping to consider when the word "otherwise" appears and asking if a better way is available.
Yes, the question came from the new Maths in Focus Ext 2 textbook in ch 5 Mathematical Induction Exercise 5.04 Q3 part b.
#### Pedro123
##### Member
Heads up - using the method of differentiation (Which is significantly easier) and induction for such problems is quite common in a lot of 3 unit exams, so just make sure you are aware of that process | 2020-03-28T18:14:52 | {
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