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https://www.jiskha.com/questions/561464/The-length-of-a-rectangle-is-2-meters-more-than-its-width-If-the-area-is-63-square
# algebra The length of a rectangle is 2 meters more than its width. If the area is 63 square meters, what are the length and width? 1. L=w+2 L x w = 63 (w+2) x w = 63 w^2+2w=63 w^2+2w-63=0 What two factors multiply to -63 and add up to 2? 9 and -7. L=9, W=7 posted by Jen First Name ## Similar Questions 1. ### pre-algebra Can you help me figure this? The area of a rectangle is 14 square meters. Find the length and width of the rectangle if it's length is 5 meters greater than its width. Use an equation and the formula for the area of a 2. ### algebra The area of a rectangle is 21 square meters. Find the length and width of the rectangle if its length is 4 meters greater than its widith. using Area of a rectangle = (width)(Length) 3. ### algebra The length of a rectangle is 2 meters more than its width. If the area is 35 square meters, what are the length and width? 4. ### calculus 1. A rectangle has a constant area of 200 square meters and its length L is increasing at the rate of 4 meters per second. a. Find the width W at the instant the width is decreasing at the rate of .4 meters per second. b. At what 5. ### calculus 1. A rectangle has a constant area of 200 square meters and its length L is increasing at the rate of 4 meters per second. a. Find the width W at the instant the width is decreasing at the rate of .4 meters per second. b. At what 6. ### algebra the length of a rectangle is 5 times its width. if the length is decreased by 3 meters, and the width is increased by 10 meters, the perimeter will be 374 meters. find the length of the original rectangle. 7. ### Math A rectangle is 40 meters longer than its width if its length were increased by 10 meters and its width would be decreased by 5 meters its area would be increased by 190 square meters. Find the dimensions 8. ### basic geometric 10. a rectangle has width the same as a side of a square whose perimeter is 20m. the length of the rectangle is 9m. find the perimeter of this rectangle. 34. The width of a rectangular picture is one-half the length. The perimeter 9. ### Maths Initially, school field is 50 meters in length with 30 meters width. After expansion, the length increased by 20% and the area becomes 2340 square meters, therefore the width has increased_________ meters. 10. ### math The length of a rectangle is 2 meters more than its width. If the area is 35 square meters, what are the length and width? (Points : 4) More Similar Questions
2018-08-16T08:55:04
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http://math.stackexchange.com/questions/739719/help-on-finding-eigenvalues-of-transformation-on-matrices
# Help on finding eigenvalues of transformation on matrices T is linear transformation working on 2x2 matrices: T(A) = $\begin{bmatrix}1 & 1\\1 &1\end{bmatrix}$ A as far as I see only 0 is an eigen value but someone told me 2 is eigen value too and I can't understand why , can anyone help please? - Because $\begin{bmatrix}1&1\\1&1\end{bmatrix}\begin{bmatrix}1\\1\end{bmatrix} = 2 \begin{bmatrix}1\\1\end{bmatrix}$. – Algebraic Pavel Apr 4 '14 at 16:19 $$T\left(\begin{bmatrix}1&1\\1&1\end{bmatrix}\right)=\begin{bmatrix}1&1\\1&1\end{bmatrix}\begin{bmatrix}1&1\\1&1\end{bmatrix}=\begin{bmatrix}2&2\\2&2\end{bmatrix}=2\begin{bmatrix}1&1\\1&1\end{bmatrix}$$ - Using the general method to find eigenvalues (and the characteristic polynomial of a matrix, of course): $$\det(\lambda I-A)=\begin{vmatrix}\lambda-1&-1\\ -1&\lambda-1\end{vmatrix}=\lambda^2-2\lambda=\lambda(\lambda-2)=0\iff \lambda=\begin{cases}0\\2\end{cases}$$ - Since $T$ is a linear map $\mathbb R^4\to\mathbb R^4$ shouldn't we be looking at a polynomial of degree $4$? – Christoph Apr 4 '14 at 16:41 Probably perhaps of the poor formating but I didn't see that: I thought it was a map on two dimensional space. Yes, you're right, yet the basics remain the same. – DonAntonio Apr 4 '14 at 16:43 In any event the problem reduces to the ordinary matrix-vector case $T \vec v = \lambda \vec v$; see my answer. – Robert Lewis Apr 4 '14 at 20:58 Well, $0$ is an eginvalue since $T(\begin{bmatrix} 1 & 1 \\ -1 & -1 \end{bmatrix}) = 0, \tag{1}$ as is most easily seen. Likewise we have $T(\begin{bmatrix} 1 & 1 \\ 1 & 1 \end{bmatrix}) = \begin{bmatrix} 2 & 2 \\ 2 & 2 \end{bmatrix} = 2 \begin{bmatrix} 1 &1 \\ 1 & 1 \end{bmatrix}, \tag{2}$ so $2$ is also an eigenvalue. So far we have proceeded by intuition an (hopefully intelligent) guesswork, but we can sytematically find all the eigenvalues and eigenvectors of $T$ by observing that, if we write $A$ in columnar form, that is, $A = [\vec a_1 \; \vec a_2 ], \tag{3}$ where $\vec a_1$ and $\vec a_2$ are the columns of $A$, then $TA = [T\vec a_1 \; T \vec a_2], \tag{4}$ so that $TA = \lambda A$ becomes $TA = [T\vec a_1 \; T \vec a_2] = \lambda [\vec a_1 \; \vec a_2 ] = [\lambda \vec a_1 \; \lambda \vec a_2]; \tag{5}$ (5) shows that the "eigenmatrices" of the operator $T$ associated with eigenvalue $\lambda$ are precisely those non-zero matrices $A$ whose non-zero columns are ordinary $\lambda$-eigenvectors of $T$. Thus we may find the eigenvalues of $T$ as it operates on $2 \times 2$ matrices by simply finding the ordinary eigevalues of $T$ as it maps vectors to vectors. Thus we construct the characteristic polynomial $p_T(\lambda)$ of $T$: $p_T(\lambda) = \det (T - \lambda I) = \det(\begin{bmatrix} 1 - \lambda & 1 \\ 1 & 1 - \lambda \end{bmatrix}) = (1 - \lambda)^2 - 1 = \lambda^2 - 2\lambda; \tag{6}$ we see the roots of $p_T(\lambda)$ are $0$ and $2$, as indicated by our intuitive approach. However, (6) shows that $0, 2$ are the only eigevalues of $T$, whether it is applied to vectors or matrices. And the "eigenmatrices" may be constructed from the ordinary eigenvectors of $T$; since the $0$-eigenvector of $T$ is $\vec e_0 = (1, -1)^T$ and the $2$-eigenvector is $\vec e_2 = (1, 1)^T$, we easily see that the matrices $[\vec e_0 \; 0], [0 \; \vec e_0] \tag{7}$ are a basis for the $2$-dimensional $0$-eigenspace of $T$ as it applies here to $2 \times 2$ matrices; likewise $[\vec e_2 \; 0], [0 \; \vec e_2] \tag{8}$ form a basis for the $2$-eigenspace; our "guesswork" matrices are simply the sums of the matrices in (7) and (8) respectively; each eigenspace of $T$ as applied to matrices is two-dimensional, in a sense composed of two copies of the corresponding eigenspaces of $T$ as applied to vectors. In this sense the answer given by DonAntonio is correct in spirit, for it acknowledges, though tacitly, that the evaluation of eigenvalues in fact reduces to, and is determined by, the action of $T$ on simple vectors. Hope this helps. Cheerio, and as always, Fiat Lux!!! -
2016-02-10T22:19:57
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https://tutorme.com/tutors/4671/interview/
TutorMe homepage Subjects PRICING COURSES Start Free Trial Julianne L. Mathematics major, Education minor at Binghamton University Tutor Satisfaction Guarantee Linear Algebra TutorMe Question: Let A and v be the matrix and vector given by: A=|0 0 2| |-3 1 6| |0 0 1 | v= |1| |3| |0| Find Av. What does this say about eigenvectors and eigenvalues of A? Find the characteristic polynomial of A. Find all eigenvalues of A. Julianne L. Av is the product of A and v. It is a 3x1 matrix where the entry in each row is the product of the corresponding row of A dot multiplied with the column on v. The first row is (0)(1)+(0)(3)+(2)(0)=0. The second row is (-3)(1)+(1)(3)+(6)(0)=0. The third row is (0)(1)+(0)(3)+(1)(0)=0. Therefore Av= | 0 0 0 | Because Av=0v, v is an eigenvector of A with eigenvalue 0 by the definitions of eigenvector and eigenvalue The characteristic polynomial of A is the determinant of the matrix (A-MI) I will be using M in place of the commonly used lambda. This matrix is | -M 0 2 | | -3 1-M 6 | | 0 0 1-M | The determinant is found by taking (-1)^(1+row#)*(the first entry in the row)*det(matrix obtained by deleting that row and the first column) for each row and adding these values. So the determinant is -M*det |1-M 6 | + 0 + 2*det |-3 1-M | | 0 1-M | | 0 0 | =-M *((1-M)*(1-M)) = -M (1-2M+M^2) = -M+2M^2-M^3 = -M*(1-M)^2 Any of these are proper expressions of the characteristic polynomial. The eigenvalues are the roots of the characteristic polynomial so all you need to do to find them is to set -M*(1-M)^2=0 and solve for M. We find that in this case M=0 and M=1. These are the eigenvalues of A. Calculus TutorMe Question: Find two linearly independent power series solutions of the differential equation y''-xy'-2y=0 Write the answer to at least x^5. Julianne L. We will be substituting in y=(summation n=0->inf) ((Cn)*x^n) Cn is a coefficient C subscript n By differentiating we find that y'=(summation n=1->inf)( (Cn)*n*x^(n-1) ) and that y''=(summation n=2->inf) ( (Cn)*n*(n-1)*x^(n-2) ) By substituting these into the original equation, our new equation is: (sum n=2->inf)((Cn)*n*(n-1)*x^(n-2)) -x*(sum n=1->inf)((Cn)*n*x^(n-1)) -2*(sum n=0->inf) ((Cn)*x^n) =0 the (-x) and (-2) can be moved inside the summations: (sum n=2->inf)((Cn)*n*(n-1)*x^(n-2)) +(sum n=1->inf)(-(Cn)*n*x^n) +(sum n=0->inf) (-2(Cn)*x^n)=0 We want to add the summations together so we will reindex the first using k=n-2 and the second two by k=n so that all of the summations contain x^k: (sum k=0->inf)((C(k+2))(k+2)(k+1)x^k) +(sum k=1->inf)(-(Ck)*k*x^k) +(sum k=0->inf) (-2(Ck)*x^k) =0 Now if we take out the k=0 term for the first and third summations, we can add all of the summations into one: 2*C2-2*C0+ (sum k=1->inf)( [(C(K+2))(k+2)(k+1)-(Ck)*k-2(Ck)]x^k ) =0 Using 2C2-2C0=0 We find that C0=C2 must be true Now we plug in values for k starting at 1 to find equations for the other coefficients, I will skip showing the computations for these as it is just plugging in integers for the Ks within the summation. Using k=1 we find that C3=(1/2)C1 with k=2 C4=(1/3)C2 and k=3 C5=(1/4)C3 Now all we need to do is pick a C0 and C1 for our two power series y1 and y2 For y1 we will choose C0=0 and C1=1. Then C2=0 C3=1/2 C4=0 C5=1/8 For y2 choose C0=1 and C1=0. Then C2=1 C3=0 C4=1/3 C5=0 Finally, we can put together our solutions. The Cn are the coefficients for each x^n term so y1= x + (1/2)x^3 + (1/8)x^5 and y2=1 + x^2 + (1/3)x^4 Calculus TutorMe Question: What is the triple integral of e^z over E, the space inclosed by the paraboloid z=1+x^2+y^2, the cylinder x^2+y^2=25, and the (x,y) plane. Use cylindrical coordinates. Julianne L. By subtracting (25=x^2+y^2) from (z=1+x^2+y^2) and simplifying, we find that z=24 is where the cylinder and paraboloid overlap. Drawing a picture of the space E shows that the paraboloid lies above the cylinder under z=24, therefore the top bound for z is the paraboloid and the bottom bound for z is the (x,y) plane or z=0. To find the triple integral using cylindrical coordinates, the limits must be in terms of z, r, and theta. The lower limit for z is z=0 and the upper limit is z=1+x^2+y^2 which is equivalent to z=1+r^2 in cylindrical coordinates since r^2=x^2+y^2. The largest radius is the radius of the cylinder x^2+y^2=25 which is 5, therefore the radius spans from 0 to 5. Theta spans from 0 to 2(pi) since the cylinder and paraboloid span a full 2(pi) radians. You additionally need to multiply an "r" within the integral since it is the Jacobian of cylindrical coordinates. The integral is then (0->2(pi)) (0->5) (0->(1+r^2)) r*e^z dz dr d(theta) Beginning with the innermost integral, r*e^z integrated with respect to z is still r*e^z and after substituting in the limits, this comes out to be (r*e^(1+r^2)-r) (r*e^(1+r^2)-r) integrated with respect to r is (1/2)e^(1+r^2)-(1/2)r^2 and this becomes (1/2)e^26-(1/2)e-(25/2) after substituting in the limits (1/2)e^26-(1/2)e-(25/2) integrated with respect to theta simply adds a theta to each element and substituting in the limits give us the final answer: (e^26)(pi)-25(pi)-e(pi) Send a message explaining your needs and Julianne will reply soon. Contact Julianne
2018-10-23T20:57:21
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https://math.stackexchange.com/questions/1268593/is-homology-determined-by-cohomology
Is homology determined by cohomology? I am aware of the universal coefficients theorem for cohomology which implies that the homology groups completely determine the cohomology groups. I am wondering if cohomology determines homology in a similar sense? If two spaces have the same cohomology groups do they necessarily have the same homology groups? Can we compute the homology groups from the cohomology groups? I am aware of Poincare duality - this only applies to "nice" spaces - seeing as homology determines cohomology in the above sense for general spaces, I was wondering if such a fact still holds switching the roles of homology and cohomology. The universal coefficients theorem for homology tells us about how the homology group with arbitrary coefficients relates to the integer homology groups - so I suppose that this is not the result I am looking for. Thanks! • It's a standard fact in algebraic topology that for closed oriented $n$-manifolds $M$, $H_{n - k}(M) \cong H^k(M)$. Look up Poincare duality. – Balarka Sen May 5 '15 at 17:51 • @BalarkaSen How is that relevant? The post doesn't even contain the word "manifold." – Matt Samuel May 5 '15 at 17:54 • @MattSamuel It's relevant in the sense that one at least knows that for (good) manifolds, two spaces with same homology has the same cohomology and vice versa, which is what the question asks. – Balarka Sen May 5 '15 at 17:55 Suppose that $X$ is a space such that all of its homology groups are finitely generated. This holds, for example, if $X$ is a "levelwise finite" CW complex (finitely many cells in each dimension), which is a very broad class of spaces. Then universal coefficients implies that all of the cohomology groups are also finitely generated. Moreover, the isomorphism class of each homology and cohomology group is determined by rank and torsion subgroup, and universal coefficients implies that • $H^k(X, \mathbb{Z})$ and $H_k(X, \mathbb{Z})$ have the same rank, and • The torsion subgroup of $H^k(X, \mathbb{Z})$ is isomorphic to the torsion subgroup of $H_{k-1}(X, \mathbb{Z})$. Hence the two sequences of isomorphism classes of groups determine each other in this case. They are nearly the same, except that the torsion subgroups are shifted one degree. Without the finitely generated hypothesis this is no longer true. For every abelian group $A$ and positive integer $n \ge 1$ it is possible to construct a Moore space $X = M(A, n)$, which is a space whose homology vanishes except in degree $n$, where it is isomorphic to $A$, and in degree $0$, where it is $\mathbb{Z}$. By universal coefficients, the cohomology of a Moore space is • $H^0(X, \mathbb{Z}) \cong \mathbb{Z}$ • $H^n(X, \mathbb{Z}) \cong \text{Hom}(A, \mathbb{Z})$ • $H^{n+1}(X, \mathbb{Z}) \cong \text{Ext}^1(A, \mathbb{Z})$ and all other cohomology vanishes. So the question is whether an abelian group $A$ is determined up to isomorphism by the isomorphism class of $\text{Hom}(A, \mathbb{Z})$ and $\text{Ext}^1(A, \mathbb{Z})$, and the answer is no. Counterexamples cannot be finitely generated: the first one that comes to mind is the following. We have $$\text{Hom}(\mathbb{Q}, \mathbb{Z}) = 0$$ (straightforward) and $$\text{Ext}^1(\mathbb{Q}, \mathbb{Z}) \cong \mathbb{R}$$ (nontrivial), and $\text{Ext}^1(-, -)$ preserves direct sums in the first factor, so it follows that the groups $\mathbb{Q}$ and $\mathbb{Q} \oplus \mathbb{Q}$ cannot be distinguished this way, since the two groups are not isomorphic, but $\mathbb{R} \cong \mathbb{R} \oplus \mathbb{R}$ (as abelian groups). • Do you know an example (Moore space or otherwise) which has trivial cohomology but non-trivial homology? – Michael Albanese Apr 7 '16 at 1:02 • @Michael: by the universal coefficient theorem this is equivalent to finding an abelian group $A$ such that $\text{Hom}(A, \mathbb{Z})$ and $\text{Ext}^1(A, \mathbb{Z})$ are both trivial (so an example exists iff an example exists which is a Moore space), assuming that you mean reduced cohomology. I don't know an example off the top of my head and am not sure if one exists. You could probably ask this as a new question. – Qiaochu Yuan Apr 7 '16 at 3:39 • For those interested, I asked this as a new question. – Michael Albanese Apr 9 '16 at 19:38
2020-04-02T10:58:26
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https://www.jakesonline.org/bear-stearns-pzzvpo/mixed-fraction-calculator-28e249
A mixed fraction is a fraction of the form \$$c {n \over d}\$$, where \$$c\$$ is an integer and \$$n d\$$. Convert the fraction to a mixed number by using long division to find the quotient and remainder 5 ÷ 3 = 1 R2 The quotient will be the whole number in the fraction, and the remainder will be the numerator in the mixed fraction 5 3 = 1 2 3 You can use … Multiply Mixed Numbers Calculator. We also offer step by step solutions. Formulas. Equivalent decimals (D) and reduced Fractions (R) will appear underneath. Study of mathematics online. Here is an adding and subtracting mixed numbers calculator to find the addition and subtraction of mixed fractions. Whether it is addition, subtraction, multiplication, or division. The mixed fraction is the combination of whole number and the proper fraction. Online Calculators for Operations with Simple and Mixed Fractions, including comparison. It is therefore the sum of a whole number and a proper fraction. A mixed fraction is a whole number and a proper fraction combined, i.e. The two fractions … DOSSIER.NET sommario. See how each example is made up of a whole number anda proper fraction together? First select if you want to use the default or mixed fraction calculator. The procedure to use the adding mixed fractions calculator … Fraction Calculators: Mixed Fraction Calculator, Mixed Number Calculator, Adding Fractions Calculator, Multiplying Fractions Calculator, Dividing Fractions Calculator, and Subtracting Fractions Calculator. One and three-half should be typed as 1 3/2. Dividing fractions calculator. Make your calculation easy by this mixed fraction calculator. The calculator performs basic and advanced operations with mixed numbers, fractions, integers, decimals. The calculator provided returns fraction inputs in both improper fraction form, as well as mixed number form. For example, to convert 5 4/5, we will do the following. Supports evaluation of mixed fractions (e.g. Yes, with this fraction calculating tool you can operate with both mixed and … An improper fraction is a fraction where the numerator (top number) is larger than the denominator (bottom number). To divide two fractions, MiroCalc.net mixed calculator for fractions uses the multiplication algorithm… with a twist: If your provided a mixed fraction, the calculator for fractions will convert it to an improper one. It is composed by the Numerator placed above a dash (called Fraction … You just need to enter your two fractions, then select the operation you want to perform. Mixed Numbers Calculator (also referred to as Mixed Fractions): This online calculator handles simple operations on whole numbers, integers, mixed numbers, fractions and improper fractions by adding, subtracting, dividing or multiplying. The mixed fraction calculator simplifies the result and displays it. Fraction Calculator. Learn to add, subtract, multiply and divide fractions. Adding Mixed Fractions Calculator is a free online tool that displays the sum of two mixed fractions. Calculate: enter 2 or 3 Fractions, select arithmetic operators using drop-downs and click on [=] button to get result. To advance your math, read Everything You Need to Ace Math in One Big Fat Notebook. Visual Fractions. Step 2) We add, subtract, multiply, or divide those improper fractions together from Step 1. The improper fraction can be converted to the mixed fraction. Keep exactly one space between the whole number and fraction and use a forward slash to input fractions. The primary reason is that the code utilizes Euclid's Theorem for reducing fractions which can be found at How MiroCalc.net Fractions Calculator is Dividing Fractions. Mixed Fraction to Improper Fraction Calculator is a free online tool that displays the conversion of mixed fraction into an improper fraction. It accepts proper, improper, mixed fractions and whole number inputs. For example, \$${5 \over 4}\$$. "Simple and Mixed". Fractions Calculator A simple, but powerful, fraction calculator to solve multiple types of fraction math problems. It is therefore the sum of a whole number and a proper fraction. A mixed fraction is a whole number and a proper fraction combined, i.e. Converting between fractions and decimals: Converting from decimals to fractions … An __improper fraction__ is a fraction whose nominator is greater than its denominator. Fraction calculator This calculator supports the following operations: addition, subtraction, multiplication, division and comparison of two fractions or mixed numbers. one and three-quarters. Electrical Calculators Real Estate Calculators When it will convert fraction to percent, a step by step procedure, results in percentage and result in decimal will be displayed separately. The combination of a proper fraction and whole number is called as mixed fraction. This unique tool will simplify your mixed number to its lowest form. © 2006 -2020CalculatorSoup® There general steps to subtract fractions are described below. This online calculator handles simple operations on whole numbers, integers, mixed numbers, fractions and improper fractions by adding, subtracting, dividing or multiplying. You can enter up to 3 digits in length for each the numerators and denominators (e.g., 456/789). Unless you got the gift of scientific mind, you need calculator to make sure you will get the correct computation during the mathematical operation with mixed fractions. Simplify Mixed Numbers. Example: Convert the following mixed number to an improper fraction. All rights reserved. Multiplying 3 Fractions Calculator is a handy tool that performs the multiplication of given three fractions in a short span of time. Sign in Log in Log out Use the online mixed fraction calculator to practice calculations with mixed fractions. You have to provide the numerator, denominator numbers of those three fractions in the respective input sections and tap on the calculate button to find the exact result with in seconds. So, in this case, you will need to divide the fractions: 1 (3/5) and 1 (4/7). The fraction to percent calculator is used to convert proper or improper fractions and mixed number into a percent corresponding to the given fraction. Do math calculations with mixed numbers (mixed fractions) performing operations on fractions, whole numbers, integers, mixed numbers, mixed fractions and improper fractions. A mixed number is a combination of a whole number and a fraction. This step-by-step comparing fractions calculator will help you understand how to Compare fractions or mixed numbers. It accepts proper, improper, mixed fractions and whole number inputs. Step 1) Convert mixed number to an improper fraction. This calculator for simply mixed fractions and allows you to change a mixed number to an improper fraction/proper fraction or vice versa. Mixed number refers to any whole number next to a fraction, for example 1 2/4. Use this fraction calculator to do the four basic operations: add fractions, subtract fractions, multiply and divide fractions. The slash separates the numerator (number above a fraction line) and denominator (number below). You can enter up to 3 digits in length for each whole number, numerator or denominator (123 456/789). If you want to simplify an individual fraction into lowest terms use our It also shows detailed step-by-step information about the fraction calculation procedure. Please, input values in this format: a b/c or b/c. Use this fractions calculator to easily perform calculations with fractions. BYJU’S online mixed fraction to improper fraction calculator tool performs the calculation faster and it displays the conversion value in a fraction of seconds. It can also convert mixed numbers to fractions and vice versa. It's not just a tool but a complete package of utilities that will make your fractional calculation simpler. It can also convert mixed numbers to fractions and vice versa. Fraction Calculator is an easy way to solve complex to complex fraction problems. BYJU’S online mixed fraction to improper fraction calculator tool performs the calculation faster and it displays the conversion value in a fraction of seconds. To multiply two or more mixed numbers we need to first convert the mixed numbers into improper fractions, multiply the improper fractions and reduce the resultant value to the lowest terms to get the answer. How to Use the Adding Mixed Fractions Calculator? Read our. For simplifying, our mixed number simplifier calculator converts the mixed number into the improper fraction. Here we have presented you a mixed fraction calculator. Sign in Log in Log out About. With this free App you can solve fraction related problems and exercices such as: Addition, Subtraction, Multiplication, Division, Simplification, Comparasion between Proper, Improper and Mixed Fractions. Example (Click to view) 1 1/3 + 2 1/4 Try it now. Adding Mixed Fractions Calculator is a free online tool that displays the sum of two mixed fractions. Use this calculator to convert your improper fraction to a mixed fraction. Fractions - use the slash “/” between the numerator and denominator, i.e., for five-hundredths, enter 5/100. The fraction calculator will generate a step-by-step explanation on how to obtain the results in the REDUCED FORM! TOGGLE 2-OR 3. Convert the fraction part of the mixed number to a decimal 2. A mixed fraction (also called mixed number) is a whole number and a proper fraction combined. Cite this content, page or calculator as: Furey, Edward "Mixed Numbers Calculator"; CalculatorSoup, 1/2 + 1/3 = (1×3+1×2) / (2×3) = 5/6 No rights can be derived from the information and calculations on this website. Use the algebraic formula for addition of fractions: Reduce fractions and simplify if possible, Use the algebraic formula for subtraction of fractions: a/b - c/d = (ad - bc) / bd, Use the algebraic formula for multiplying of fractions: a/b * c/d = ac / bd, Use the algebraic formula for division of fractions: a/b ÷ c/d = ad / bc. Fraction format button is used to work with all fractions. Enter fractions and press the = button. Simplify Online Fractions. That is why it is called a "mixed" fraction (or mixed number). 3 FRACTIONS, DECIMALS AND MIXED NUMBERS. Online Comparing Fractions Calculator… Mixed fractions are also called mixed numbers. The answer is provided in a reduced fraction and a mixed number if it exists. If you are using mixed numbers, be sure to leave a single space between the whole and fraction part. Mixed numbers: Enter as 1 1/2 which is one and one half or 25 3/32 which is twenty five and three thirty seconds. Compare: subtract second Fraction … Practical Example #1: Let’s say that you are looking to subtract: 4 (1/8) – 1 (1/5) So, as we mentioned, the first thing that you need to do is to make sure that you have the mixed fractions selected and not the regular fractions. Look for a button that has a black box over a white box, x/y, or b/c. For example, \$${5 \over 4}\$$. Push this button to open the fraction feature on your calculator. $$\dfrac{a}{b} + \dfrac{c}{d} = \dfrac{(a \times d) + (b \times c)}{b \times d}$$, $$1 \dfrac{2}{6} + 2 \dfrac{1}{4} = \dfrac{8}{6} + \dfrac{9}{4}$$, $$= \dfrac{(8 \times 4) + (9 \times 6)}{6 \times 4}$$, $$= \dfrac{32 + 54}{24} = \dfrac{86}{24} = \dfrac{43}{12}$$, $$\dfrac{a}{b} - \dfrac{c}{d} = \dfrac{(a \times d) - (b \times c)}{b \times d}$$, $$\dfrac{a}{b} \times \dfrac{c}{d} = \dfrac{a \times c}{b \times d}$$, $$\dfrac{a}{b} \div \dfrac{c}{d} = \dfrac{a \times d}{b \times c}$$. Greatest Common Factor Calculator. With this online mixed fraction calculator (or mixed number calculator) with whole numbers and fractions you can easily add mixed fractions, subtract mixed fractions, multiply mixed fractions and divide mixed fractions. With this calculator you can practice calculations with mixed fractions (mixed numbers). Dividing fractions calculator online. Fraction Operations and Manipulations Click on the decimal format button, enter a fraction or mixed number, then click equals. A mixed number is a combination of a whole number and a fraction. This calculator divides two fractions. … BYJU’S online adding mixed fractions calculator tool makes the calculation faster and it displays the sum of two mixed fractions in a fraction of seconds. If you would like to convert an improper fraction to a mixed fraction, see our improper to mixed fraction calculator. Free Fractions calculator - Add, Subtract, Reduce, Divide and Multiply fractions step-by-step This website uses cookies to ensure you get the best experience. Adding and subtracting mixed number is confusing. This calculator allows you to add, subtract, multiply, or divide simple fractions and mixed fractions online. This video will show you how to convert from mixed to improper fractions as well as improper to mixed fractions using the calculator Casio fx-991ms Mixed number to fraction conversion calculator that shows work to represent the mixed number in impropoer fraction. Fractions Calculator. Mixed fractions are also called mixed numbers. Add, Subtract, Multiply, Divide, Compare Mix or Compound Fractions … Simplify Fractions Calculator. If you are simplifying large fractions by hand you can use the If they exist, the solutions and answers are provided in simplified, mixed and whole formats. Order. The fraction calculator is very helpful in real life. This online fraction calculator will help you understand how to add, subtract, multiply or divide fractions, mixed numbers (mixed fractions), whole numbers and decimals.The fraction calculator will compute the sums, differences, products or quotients of fractions. A mixed fraction is a fraction of the form c n d, where c is an integer and n < d. For example, 11 4 = 2 3 4. Mixed fraction … This mixed fraction calculator by calculator-online is the smart tool that helps you in adding, subtracting, multiplying, and dividing mixed numbers fraction. You have 6 sections: 1. The Fraction Calculator will reduce a fraction to its simplest form. If the fraction or mixed number is only part of the calculation then omit clicking equals and continue with the calculation per usual. You can also add, subtract, multiply, and divide fractions, as well as, convert to a decimal and work with mixed numbers and reciprocals. Study math with us and make sure that "Mathematics is easy!" Fill in two fractions and choose if you want to add, subtract, multiply or divide and click the "Calculate" button. There general steps to divide fractions are described below. 1/2 + 1/3 = (1×3+1×2) / (2×3) = 5/6 The answer will be in fraction form and whole/decimal form. Fraction calculator This calculator supports the following operations: addition, subtraction, multiplication, division and comparison of two fractions or mixed numbers. Advatages: It is easy to use. You can also add, subtract, multiply, and divide fractions, as well as, convert to a decimal and work with mixed numbers and reciprocals. Mixed Fraction 1 : Operator : Mixed Fraction 2 . Read on to learn about fraction percent and its formula and how to turn a fraction … Use this calculator to convert your improper fraction to a mixed fraction. Adding mixed numbers, converting fraction to whole number, multiplying fractions by whole numbers, subtracting mixed numbers, and multiplying mixed fractions are among the processes this calculator can do. Enter the value of two mixed numbers and click calculate, it displays the step by step solution with the answer. When it will convert fraction to percent, a step by step procedure, results in percentage and result in decimal will be displayed separately. The Fraction Calculator will reduce a fraction to its simplest form. Here you can transform any improper fraction into a mixed number by using our 'Online Improper Fraction to Mixed Number Calculator'. Open the fraction calculator to practice calculations with mixed numbers calculator to solve types... Lowest forms by dividing both numerator and denominator, i.e., for,. Free online tool that displays the conversion of mixed fraction is a number! Would like to convert an improper fraction to power ( fraction or mixed numbers improper! Return the result and displays it the numerator ( top number ) is larger than the denominator ( 123 )! Into a percent corresponding to the mixed fraction… the fraction feature is on, should... Faster than others you might find number simplifier calculator converts the mixed fraction is a,... 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2021-05-17T01:48:57
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http://mathhelpforum.com/algebra/225196-watch-gaining-losing-time.html
# Math Help - Watch gaining and losing time 1. ## Watch gaining and losing time My watch (which is a 12 hour watch) GAINS 3 mins every 2 hours. a) I set my watch to the correct time at noon on 1st January. If I don't reset it, when will it next show the correct time? I got: real time vs gain time 2hrs vs 3 mins gain 40hrs vs 60 mins (x20) 960 hrs vs 24 hrs (x24) = 40 days (divided 960 hrs by 24hrs) Is this correct please for a? b) Mrs Warma's watch (also 1 12 hour watch LOSES 5 mins every 2 hours. She also sets her watch to the correct time at noon on 1st January. When will our two watches next show the same time? I got: real time vs lose time 2hrs vs 5 mins lose I am stuck after this.... c) When will our watches next show the same, correct time? I am stuck with this 2. ## Re: Watch gaining and losing time Originally Posted by Natasha1 My watch (which is a 12 hour watch) GAINS 3 mins every 2 hours. a) I set my watch to the correct time at noon on 1st January. If I don't reset it, when will it next show the correct time? I got: real time vs gain time 2hrs vs 3 mins gain 40hrs vs 60 mins (x20) 960 hrs vs 24 hrs (x24) = 40 days (divided 960 hrs by 24hrs) Is this correct please for a? b) Mrs Warma's watch (also 1 12 hour watch LOSES 5 mins every 2 hours. She also sets her watch to the correct time at noon on 1st January. When will our two watches next show the same time? I got: real time vs lose time 2hrs vs 5 mins lose I am stuck after this.... c) When will our watches next show the same, correct time? I am stuck with this Not quite. It's a 12 hr watch so it will be correct again in 12 hrs, not 24 as you've used. So it will be 20 days not 40. b) $\mbox{in minutes your watch shows }m1=t(1+3/120)=t(1+1/40)=\frac{41}{40}t\mbox{ minutes.}$ $\mbox{her's shows }m2=t(1-5/120)=t(1-1/24)=\frac{23}{24}t\mbox{ minutes.}$ $\mbox{The watches are able to display } t \bmod 720 \mbox{ minutes uniquely.}$ $\mbox{So we want to solve for }\left(m1 \bmod 720\right)=\left(m2 \bmod 720\right)$ see if you can work it from there. c) you need to calculate how often the lady's watch is correct similar to how you did in (a). Then find the least common multiple of the three times, how often your watch is right, how often hers is, how often they are the same. 3. ## Re: Watch gaining and losing time i don't understand what mod is, never done this at school 4. ## Re: Watch gaining and losing time a mod b = remainder of a/b for example (25 mod 7) = 4 because 25 - 3x7 = 4 I'll give you a hint solving this. Try solving $\frac{41}{40}t-720=\frac{23}{24}t$ 5. ## Re: Watch gaining and losing time Hello, Natasha1! My watch (which is a 12-hour watch) GAINS 3 mins every 2 hours. a) I set my watch to the correct time at noon on 1st January. If I don't reset it, when will it next show the correct time? Your watch gains 1.5 minutes every hour. To show the correct time again, it must gain 12 hours = 720 minutes. . . $\frac{720\text{ min.}}{1} \times \frac{1\text{ hr.}}{1.5\text{ min.}} \:=\:480\text{ hours} \:=\:20\text{ days}$ It will show the correct time at noon on 21 January. 6. ## Re: Watch gaining and losing time Originally Posted by Natasha1 b) Mrs Warma's watch (also 1 12 hour watch LOSES 5 mins every 2 hours. She also sets her watch to the correct time at noon on 1st January. When will our two watches next show the same time? Your two watches will show the same time at the moment when your watch gains $t$ minutes and Mrs. Warma's watch loses 12 hours minus $t$ minutes. So, let $h$ be the number of periods of two hours that must pass before your watch gains $t$ minutes (and Mrs. Warma's watch loses $720-t$ minutes). Then $3h = t$ and $5h = 720-t$. This means $t=3h$ and $t=720-5h$. So, $3h = 720-5h$ shows that $8h = 720$, and $h = 90$. That means their watches will show the same time every $2h = 180$ hours. So, their watches will show the same time midnight, January 9th (that is the midnight that falls between the 8th and the 9th). 7. ## Re: Watch gaining and losing time Hello again, Natasha1! b) Mrs Warma's watch (a 12-hour watch) LOSES 5 mins every 2 hours. She sets her watch to the correct time at noon on 1st January. When will our two watches next show the same time? Your watch gains 1.5 minutes every hour. To show the correct time again, your watch must gain 12 hours = 720 minutes. . . $\frac{720\text{ min.}}{1} \times \frac{1\text{ hr.}}{1.5\text{ min.}} \:=\:480\text{ hours} \:=\:20\text{ days}$ Her watch loses 2.5 minutes every hour. To show the correct time again, her watch must lose 12 hours = 720 minutes. . . $\frac{720\text{ min.}}{1} \times \frac{1\text{ hr.}}{2.5\text{ min.}} \:=\:288\text{ hours} \:=\:12\text{ days}$ The two watches will agree in: $LCM(20,12) = 60$ days . . . Noon on March 30. 8. ## Re: Watch gaining and losing time Originally Posted by Soroban Hello again, Natasha1! Your watch gains 1.5 minutes every hour. To show the correct time again, your watch must gain 12 hours = 720 minutes. . . $\frac{720\text{ min.}}{1} \times \frac{1\text{ hr.}}{1.5\text{ min.}} \:=\:480\text{ hours} \:=\:20\text{ days}$ Her watch loses 2.5 minutes every hour. To show the correct time again, her watch must lose 12 hours = 720 minutes. . . $\frac{720\text{ min.}}{1} \times \frac{1\text{ hr.}}{2.5\text{ min.}} \:=\:288\text{ hours} \:=\:12\text{ days}$ The two watches will agree in: $LCM(20,12) = 60$ days . . . Noon on March 30. This is the answer to part (c), not part (b). 9. ## Re: Watch gaining and losing time Many thanks to all of you . Merry Christmas everyone!
2014-10-21T09:56:54
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https://math.stackexchange.com/questions/1027539/what-does-this-summation-mean
# what does this summation mean? I apologize if this question has been asked. I know several similar ones have been asked but I cannot find one answering this in particular. I want to know what this summation means: $$\sum_{i,j=1}^{M}$$ Is this equivalent to $$\sum_{i=1}^M \sum_{j=1}^M$$ or does it mean that i and j will be incremented together at the same time? If context is needed, the equation is coming from here equation 16 on the third page (page 392 of the journal). • Yes, they are equivalent. Nov 18 '14 at 13:05 • It's the sum over the set of indices $(i,j)$ with $1\leqslant i,j\leqslant M$ which is equivalent to $\sum_{i=1}^M\sum_{j=1}^M$. Nov 18 '14 at 13:05 • $\sum_{i,j=1}^{M}\equiv \sum_{i=1}^M \sum_{j=1}^M$ with no doubt. Nov 18 '14 at 13:06 $$\sum_{i,j=1}^{M}a_{i,j}=\sum_{1\le i,j\le M}a_{i,j}=\sum_{i=1}^M\sum_{j=1}^Ma_{i,j}=\sum_{j=1}^M\sum_{i=1}^Ma_{i,j}$$ means that we sum the terms $a_{i,j}$ with all the possible values of $i$ and $j$ between $1$ and $M$.
2021-10-16T22:15:52
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https://math.stackexchange.com/questions/3632614/r-wedge-neg-s-rightarrow-neg-q/3632666
# $(r \wedge \neg s) \rightarrow \neg q$ How can this English sentence be translated into a logical expression? In Kenneth Rosan, the answer to this following sentence “You cannot ride the roller coaster if you are under 4 feet tall unless you are older than 16 years old.” is given as, $$(r \wedge \neg s) \rightarrow \neg q$$ Where, q: “You can ride the roller coaster.” r: “You are under 4 feet tall.” s: “ You are older than 16 years old.” My solution: So, I broke down this compound sentence as follows: “[You cannot ride the roller coaster] if [you are under 4 feet tall] unless [you are older than 16 years old.]” Now, substituting variables in given compound sentence. ($$\neg q$$) if (r unless s). Applying equivalence formula for Q if P $$\Leftrightarrow$$ P $$\to$$ Q (r unless s) $$\to$$ ($$\neg q$$) Now, solving for unless. So, (r unless s) $$\Leftrightarrow$$ ($$\neg s \to r$$) ref. ($$\neg s \to r$$) $$\to$$ ($$\neg q$$) Again solving for $$\to$$ (implication), we get: (s $$\lor$$ r) $$\to$$ ($$\neg q$$) So, my derivation is obviously wrong and does not match with Kennet Rosen. My Question: What mistake I did? and How to derive the given answer systematically? As noted by Jay, Kenneth Rosen interprets (r unless s) according to: $$\begin{array}{cc|c} r & s & (r \text{ unless } s) \Leftrightarrow (\neg s \to r) \\\hline 0 & 0 & 0 \\ 0 & 1 & 1 \\ 1 & 0 & 1 \\ 1 & 1 & 1 \\ \end{array}$$ The issue turns out to be the order of operations for "unless". You started by breaking it up like this “[You cannot ride the roller coaster] if [you are under 4 feet tall] unless [you are older than 16 years old.]” And substituted variables to get: ($$\neg q$$) if ($$r$$ unless $$s$$). If instead we use a different order of operations to group these, it works out with your definition of unless. That is, we have: $$((\neg q) \mathbf{\text{ if }} r) \mathbf{\text{ unless }} s$$ Now using $$(P \mathbf{\text{ if }} Q) \Leftrightarrow (Q \to P)$$ $$(r \to \neg q) \mathbf{\text{ unless }} s$$ Now using your $$(P \mathbf{\text{ unless }} Q) \Leftrightarrow (\neg P \to Q)$$ $$\neg s \to (r \to \neg q)$$ expanding $$\neg s \to (\neg r \lor \neg q)$$ expanding $$s \lor (\neg r \lor \neg q)$$ This is the same as the other result. To see this, use associativity of logical or $$(s \lor \neg r) \lor \neg q$$ Then turn it into an implication $$\neg(s \lor \neg r) \to \neg q$$ Use demorgan's law $$(\neg s \land r) \to \neg q$$ (EDIT: Previously I arrived at the answer with the same order of operations, but a different interpretation of unless: $$(r \text{ unless } s) = (r \text{ and } \neg s)$$. Because "(anything) unless True = True" seriously sounds wrong to me. My interpretation of unless worked in this case, but apparently is not the correct english interpretation. Apologies.) • For this reason, we are going to translate ‘P unless Q’ with just ¬Q → P unless stated otherwise from rpi.edu – Ubi hatt Apr 19 '20 at 4:14 • Also, what formula you assumed for "r unless s"? – Ubi hatt Apr 19 '20 at 4:15 • @Ubihatt I used no formula. I filled it out how I interpret the phrase in english. So it appears this may just be a natural language ambiguity if you interpret it otherwise. – HBrown Apr 19 '20 at 4:17 • Kenneth Rosen also has formula for unless which exactly matches with formula embedded in my previous comment. Please check. Unless is a keyword. – Ubi hatt Apr 19 '20 at 4:19 • Interpreting "P unless Q" = "¬Q → P", would mean "You cannot ride the roller coaster if you are under 4 feet tall unless you are older than 16 years old." says you cannot ride the roller coaster if you are older than 16 years old, no matter what your height is. Is that really how you interpret that phrase? – HBrown Apr 19 '20 at 4:27 The mistake is to interpret the sentence as $$(\neg q)\ \textbf{if}\ (r\ \textbf{unless}\ s)$$ The correct interpretation is $$(\neg q\ \textbf{if}\ r)\ \textbf{unless}\ s$$ • You beat me to the punch =) Nicely done. – Jay Dunivin Apr 19 '20 at 5:01 I have a discrete math book by Kenneth Rosen, and here is an excerpt from the book listing the equivalent ways to express $$p \implies q$$ in English. The one you mentioned is boxed in blue. But where did you go wrong? I believe it was at this step Now, substituting variables in given compound sentence. ($$\neg q$$) if (r unless s). The word unless does not attach $$s$$ to $$r$$; it attaches $$s$$ with the proposition $$\neg q$$ if $$r$$. The reason for this is that there was a proposition established prior to unless, which was the statement, "You cannot ride the roller coaster if you are under 4 feet tall." That sentence is established as a proposition; it is the $$q$$ in the table I provided. So the correct statement using your scheme is $$\big( \, \neg q \, \textbf{ if } \, r \, \big) \, \textbf{ unless } \, s,$$ which is equivalent to $$\neg s \implies (r \implies \neg q).$$ Can you take it from here? • As you saw, I initially interpreted unless differently, and just happened to get the right answer erroneously. So I learned something here too. But I'm struggling how "True unless True" evaluates to True. "The dog is angry unless the dog is fed". If the dog is angry and the dog is fed, I would interpret that as false. Is my english screwed up, or what am I missing? – HBrown Apr 19 '20 at 5:43 • Maybe what you're missing is that $F \implies T$ is true (similarly, $F \implies F$ is true). "The dog is angry unless the dog is fed," which is the same as "If the dog is not fed, then the dog is angry." If the dog is indeed fed, then the statement "the dog is not fed" is false, of course. But then the statement "If the dog is not fed, then the dog is angry" is actually a true statement (a false hypothesis can be used to derive any conclusion). The only way the statement could be false is when the dog is not fed, but the dog is not angry. Does that help? – Jay Dunivin Apr 19 '20 at 5:49 • To add to what I said, "True unless True" is equivalent to "False implies True," which is a true statement, per the definition of "implies" that we have agreed to adopt. We have agreed that $p \implies q$ is false if and only if $p$ is true and $q$ is false, and $p \implies q$ is true otherwise. – Jay Dunivin Apr 19 '20 at 5:53 • I see where you're coming from. I don't think it's your English. I was looking at other sites and papers, and there isn't a consensus as to what "unless" means in general, for it can be construed in different ways in different contexts. This link is particularly helpful as to why in logic this particular convention was adopted: math.stackexchange.com/questions/1295106/… – Jay Dunivin Apr 19 '20 at 6:14 • You might also be interested in these other links regarding the ambiguity of "unless" outside of mathematics: Reddit: reddit.com/r/askphilosophy/comments/52585m/… and a linguistic paper: web.stanford.edu/~danlass/Nadathur-Lassiter-unless-SuB.pdf – Jay Dunivin Apr 19 '20 at 6:15
2021-06-24T13:01:09
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## Math City Derivative Formulas It has two major branches, differential calculus and integral calculus. Derivative of the Exponential Function. Advanced:(optional) Given a formula for a function obtain improved numerical estimates of the derivative by choosing smaller increments (e. These applications include acceleration and velocity in physics, population growth rates in biology, and marginal functions in economics. Polar Integral Formula: Positive Series. Computation and Properties of the Derivative in Calculus. " Be Careful: "Find the derivative using the limit definition" does not mean estimating the derivative like we did earlier. Calculus-Specific Formulas There are a number of basic formulas from calculus that you need to memorize for the exam. The golden ratio. And it doesn't just work with position; Calculus can work with any function. ACTUARIAL SCIENCE FIRST YEAR CALCULUS - FORMULA SHEET Standard Integrals R xndx = xn+1 n+1 (n 6= −1) R 1 x dx = lnx R 1 x dx= lnx R ex= R sinxdx = −cosx R cosecxcotxdx = −cosecx. You should convince yourself that both of these expressions are indeed the same, by writing out explicitly the four terms of each of the two formulas! Let's try this for our derivative formulas: We shifted each occurrence of n in the expression up by one unit, while the limits of summation were shifted down by one unit, from 1 to 0, and from to. The instantaneous rate of change measures the rate of change, or slope, of a curve at a certain instant. It was submitted to the Free Digital Textbook Initiative in California and will remain unchanged for at least two years. Table of Derivatives. – is easier than you think. We discuss the Mittag-Leffler bounds of these solutions. ) And since the derivative of a sum is the sum of “Basic integration formulas. For the first two, determine the formula for the derivative by thinking about the nature of the given function and its slope at various points; do not use the limit definition. if n and a are positive real numbers, and a is not equal to 1, then If a x = n, then log a n = x. Applications of Differentiation. MATH 19B, Integral Calulcus: Some useful Trig formulas, derivatives and integrals. Stay ahead with the world's most comprehensive technology and business learning platform. #d/dx sec(x)=sec(x)tan(x)# You could memorize this, but you can work it out too by knowing some trig properties. f(x) = (x4 +3x)−1. The Derivation Formula. $\endgroup$ – Federico Poloni Aug 17 '15 at 8:42 1 $\begingroup$ This question really belongs to math. Mara~o garded on a algebra which generalizes certain features of th e derivative operator. derivative synonyms, derivative pronunciation, derivative translation, English dictionary definition of derivative. Tutorial on elementary differentiation formulas, their derivation and use. Ask Question 3. Derivative Formulas. COMMON MATH FORMULAS MISCELLANEOUS FORMULAS Simple Interest + L L N P where I = interest, p = principal, r = interest rate, and t = time Distance @. The power rule for differentiation states that if n is a real number and f(x) = x^n, then f'(x) = nx^{n-1}. Greater than 1, there will be exponential growth towards infinity (positive). Calculus allows us to see the connection between these equations. It may be rewritten as Another similar formula is given by Derivative of the Inverse Function. Together, we are going to review why and how to generate the formula for the Definition of Derivative, and we will walk through an example of how to use this definition to. You can choose formulas from different pages. To insure your continued success in second-semester, it is important that you are able to recall and use the following facts without struggling. f(x) = ex sinx 3. Recall the definitions of the trigonometric functions. org Math Tables: Derivatives of Hyperbolics Proofs of Derivatives of Hyperbolics Proof of sinh(x) = cosh(x): From the derivative of e x. fourth, fifth), extracting more and more information from that simple position function. Here is the rewritten function. The derivative of csc x. 6 Derivatives of Inverse Functions. if n and a are positive real numbers, and a is not equal to 1, then If a x = n, then log a n = x. Derivative Formula. Present Value Formula $$\huge P = \frac{F}{(1+r)^t}$$ The present value of money is equal to the future value divided by the interest rate plus 1 raised to the t power, where t is the number of months, years, etc. Compatible numbers. A class of second derivative extended backward differentiation formulas suitable for the approximate numerical integration of stiff systems of first-order ordinary differential equations is examine. Calculus Here is a list of skills students learn in Calculus! These skills are organized into categories, and you can move your mouse over any skill name to preview the skill. as a formula for the derivative of Φ with respect to Ξ. Most of derivatives' value is based on the value of an underlying security, commodity, or other financial instrument. But calculus provides an easier, more precise way: compute the derivative. In this case, it's easy to write down explicit formulas for the basis functions (they are just Bernstein polynomials), so the algebra is simpler. The derivative of ln(x) is a well-known derivative. Then the function. And it doesn't just work with position; Calculus can work with any function. To find the slope of a line tangent to a parabola at a specific point, find the derivative of the parabola's equation, then substitute the -coordinate of the specific point in the new equation. Or, as a formula that tells you the rate of change of a certain function. A very, very, very small distance, but large enough to calculate the slope. Derivative is a rate of change of function with respect to a variable. Choose from 500 different sets of calculus 1 formulas math flashcards on Quizlet. Find the equation to the tangent line to the curve y = x+ p x at the point (1,2) Example 3. In rst-semester calculus (regardless of where you took it) you learned the basic facts and concepts of calculus. The second derivative can be used to determine the concavity and inflection point of a function as well as minimum and maximum points. One of our goals in this section is to learn how to differentiate the logarithm function. Calculus 8th Edition answers to Chapter 2 - Derivatives - 2. It may be rewritten as Another similar formula is given by Derivative of the Inverse Function. In this page you can see a list of Calculus Formulas such as integral formula, derivative formula, limits formula etc. Math skills assessment. In Topic 19 of Trigonometry, we introduced the inverse trigonometric functions. The derivative of a function f at a point x is commonly written f '(x). Explicitly: If this limit exists, then we say that the derivative exists and has this value, and we say that the function is differentiable at the point.   These questions will be similar to see how you will be tested on the exam. Firstly, you need to memorize these math formula very carefully before your exam. Calculus is the branch of mathematics studying the rate of change of quantities and the length, area and volume of objects. Includes the Power Formula. Write with me. To obtain a formula for the derivative of any expression in x, assume the function. The Derivative tells us the slope of a function at any point. Show Instructions. Offline content. This page is send by Ali Nawaz Bajwa (MS(Math), M. PHYSICS FORMULAS 2426 Electron = -1. The basic ideas are not more difficult than that. Finding the formula of the derivative function is called differentiation, and the rules for doing so form the basis of differential calculus. It is the only mathematics course certified to meet the university's "Viewing a Wider World" upper division general education requirement. Cheat Sheets and Calculator Tips The following are a list of formula sheets for various courses taught in the Mathematics Department. Please send suggestions for amendments to the Secretary of the Teaching Committee, and they will be considered for incorporation in the next edition. There are two ways of introducing this concept, the geometrical way (as the slope of a curve), and the physical way (as a rate of change). In a formula, it is abbreviated to just 'sec'. Calculus 8th Edition answers to Chapter 2 - Derivatives - 2. In general, scientists observe changing systems (dynamical systems) to obtain the rate of change of some variable. One of the most important functions in all of mathematics is the natural exponential function f(x)=ex. The derivative of ln u(). How to calculate derivatives for calculus. Offline content. When trying to gure out what to. Dick was a member of the National Academy of Sciences and a Honorary Fellow of the Indian Mathematical Society. Find the derivative of the equation and explain its physical meaning. Undergraduate. This is the original method in finding the derivative of any equations using the Increment Method. It is possible to guess at a formula for the derivative from this curve. Used by over 7,000,000 students, IXL provides personalized learning in more than 8,000 topics, covering math, language arts, science, social studies, and Spanish. It covers basic set theory and logic, relations and functions, and how to analyze, construct, and write clearly reasoned, well-structured elementary proofs using universal techniques. For each of the listed functions, determine a formula for the derivative function. You can either use an established mnemonic device or you can create your own. Calculation of the derivative — the most important operation in differential calculus. Derive the function 5/(x^0. Specific differentiation formulas You will be responsible for knowing formulas for the derivatives of these func­ tions: xn, sin−1 x, tan−1 x, sin x, cos x, tan x, sec x, ex , ln x. The derivative of sec x. Help With Your Math Homework. After you have selected all the formulas which you would like to include in cheat sheet, click the "Generate PDF" button. First derivative test for maxima/minima problems. For most of the. 331 (3/23/08) Estimating directional derivatives from level curves We could find approximate values of directional derivatives from level curves by using the techniques of the last section to estimate the x- and y-derivatives and then applying Theorem 1. The expressions are obtained in LaTeX by typing \frac{du}{dt} and \frac{d^2 u}{dx^2} respectively. We discuss the Mittag-Leffler bounds of these solutions. Calculate the Costs to Use Electricity - powered by WebMath. In fact, you can use calculus in a lot of ways and applications. However, we can generalize it for any differentiable function with a logarithmic function. Basic Properties. The derivative of y = arccsc x. Method 1 (The way I found on my own): ∞ ∑ i=1a0rn−1≡S S=a0r0+a0r1+a0r2+⋯ S=r (a0r−1+a0r0+a0r1+⋯) S=r (a0r−1+S) S=a0+rS (1−r)S=a0 S=a0 (1−r) Note that for this to work, you must first confirm this: lim n→∞an=0. Eventually, formulas are used to provide mathematical solution for real world problems. Note that we studied Exponential Functions here and Differential Equations here in earlier sections. If f(x) = (3 x )(sin x) (cos x), find f'( x ). If there is a function h, then the anti-derivative of this function will be a differential function, say H. COMMON MATH FORMULAS MISCELLANEOUS FORMULAS Simple Interest + L L N P where I = interest, p = principal, r = interest rate, and t = time Distance @. Basic math formulas Algebra word problems. Solution If you’re having trouble with this problem, it may help to review Professor. PHYSICS FORMULAS 2426 Electron = -1. Use the Fundamental Theorem of Calculus to evaluate each of the following integrals exactly. Generally, mathematical finance will derive and extend the mathematical or numerical models without necessarily establishing a link to financial theory, taking observed. Then the function. You can also send us message on facebook. Partial Derivatives Calculus III - Chapter 11 Formulas Equation of tangent plane to surface of z f(x,y) at point P(x ,y z )is: z z f (x ,y )(x x ) f (x ,y )(y y ). Learn all about derivatives and how to find them here. Finding Derivative with Definition of Derivative Calculus 1 AB - Duration: 15:27. That means there are no two x-values that have the same y-value. \) The domains of the functions and their graphs. A simple calculus check reveals that the latter is the derivative of the former with respect to R. To find the equation of the tangent line using implicit differentiation, follow three steps. In this case, the instantaneous rate is s'(2). Calculus broadly classified as Differentiation and Integration. This course is a transition from procedural mathematics, such as calculus, to advanced mathematics where proofs are the professional language of discourse. A comprehensive list of the most commonly used basic math formulas. For each of the listed functions, determine a formula for the derivative function. Functions differentiation formula In the table below u and v — are functions of the variable x , and c — is constant. Seminar on Pure Mathematics Introduction to Hamiltonian Stationary Lagrangian submanifolds and a compactness theorem in 2 dimension Room 3494, Academic Building (lift 25-26), HKUST. A list of the most commonly used algebra formulas. The First and Second Derivatives The Meaning of the First Derivative At the end of the last lecture, we knew how to differentiate any polynomial function. Learn calculus derivatives formulas math with free interactive flashcards. -- Encyclopedia Of Mathematics (Science Encyclopedia)-- The Princeton Companion to Mathematics-- Handbook of Mathematical Functions: with Formulas, Graphs, and Mathematical Tables (Dover Books on Mathematics)-- NIST Handbook of Mathematical Functions (Being of course the 2010 update of the Abramowitz classic above). 6666666666666666)) using the power rule. By applying the derivation formulas and using the usual derivation table, it is possible to calculate any function derivative. Students, teachers, parents, and everyone can find solutions to their math problems instantly. The Derivative Calculator supports computing first, second, …, fifth derivatives as well as differentiating functions with many variables (partial derivatives), implicit differentiation and calculating roots/zeros. Our study guides are available online and in book form at barnesandnoble. The derivative of acceleration is jerk. In this page you can see a list of Calculus Formulas such as integral formula, derivative formula, limits formula etc. It is fun and instructive and the student version without solutions is available free of charge. And, you need these formulas in your exams. Derivative proofs of csc(x), sec(x), and cot(x). The Derivative Measures Slope. Differentials and Newton's method. \] Now we consider the logarithmic function with arbitrary base and obtain a formula for its derivative. f(z)=−1/z^5. Calculus is used in mechanical, physics etc. 2 Functions of 2 or more variables Functions which have more than one variable arise very commonly. Once you understand the concept of a partial derivative as the rate that something is changing, calculating partial derivatives usually isn't difficult. Notation This is best described with an example. The derivative of sin x. SHOW PROPER CALCULUS-LEVEL WORK SIMPLIFY ALL ANSWERS Find lim cothx algebraically C' Hô Prove the logarithmic formula for Sinh x given in your textbook. If you are looking for a formula to solve your basic math problems, your formula is likely here. The desired formula is obtained by taking the limit of both sides as x approaches zero so that. It is designed for anyone who needs a basic to advanced understanding of mathematics concepts and operations. qxd Author: ewedzikowski Created Date: 10/29/2004 9:36:46 AM. Formulas: (a) sin2 x+cos2 x =1 (b) 1+tan2 x = sec2 x (c) cos2 x = 1+cos(2x) 2 (d) sin2 x = 1−cos(2x) 2 2. 1 Recall: ordinary derivatives If y is a function of x then dy dx is the derivative meaning the gradient (slope of the graph) or the rate of change with respect to x. Choose from 500 different sets of calculus 1 formulas math flashcards on Quizlet. Differential. Derivatives Calculus Differential Calculus Algebra Ap Calculus Calculus Notes I Love Math Fun Math Math Math Math Tutor Find all the essential math formulas you need, all in one place. This formula will be derived and then applied to … • the rôle of the Wronskian in the solution of linear differential equations,. The derivative is the function slope or slope of the tangent line at point x. Differentiate the formula with respect to time. more mathematical) definition. Every technique outlined in this article on calculating derivatives can be verified by a proper use of the definition of the derivative. 1300 Math Formulas Mega Pack - Table of Derivatives - Higher Order Derivatives - Application of Derivative Crime City 2019. Precalculus & Elements of Calculus tutorial videos. Figure 1 shows two graphs that start and end at the same points but are not the same. If y = 3 x 2, which can also be expressed as f(x)= 3 x 2, then. Derivative proofs of csc(x), sec(x), and cot(x). Taylor's Formula; List of Derivatives of Trig & Inverse Trig. And (from the diagram) we see that: Now follow these steps: Fill in this slope formula: ΔyΔx = f(x+Δx) − f(x)Δx. How to Cite This Entry: Newton-Leibniz formula. We need differentiation when the rate of change is not constant. Derivative Identities. How is Business Calculus Different?. Derivatives are fundamental to the solution of problems in calculus and differential equations. derivative synonyms, derivative pronunciation, derivative translation, English dictionary definition of derivative. We will use the combination of these modes in order to have a desired and accurate output. We simplify the equation by taking the tangent of both sides: y = tan−1 x tan y = tan(tan−1 x) tan y = x. The expression (4. Taking the derivative¶. com; [email protected] Elasticity of demand is a measure of how demand reacts to price changes. Create a chart showing the graphical relationships between f, f. com to clear your doubts from our expert teachers and download the Application of Derivatives formula to solve the problems easily to score more marks in your Board exams. 2 days ago · Universal free lunch is linked to better test scores in New York City, new report finds bumps in reading and math state test scores once students attended schools with universal free lunch. A class of computationally A-stable formulas is presented, with notes on implementing a set of them in a variable stepsize, variable order method. Total derivatives are often used in related rates problems; for example, finding the rate of change of volume when two parameters are changing with time. Find the derivative (d/dx)(x^(-0. A first-semester college calculus course and the subsequent single-variable calculus course Recommended Prerequisites You should have successfully completed courses in which you studied algebra, geometry, trigonometry, analytic geometry, and elementary functions. Implicit multiplication (5x = 5*x) is supported. 18 Useful formulas. You should convince yourself that both of these expressions are indeed the same, by writing out explicitly the four terms of each of the two formulas! Let's try this for our derivative formulas: We shifted each occurrence of n in the expression up by one unit, while the limits of summation were shifted down by one unit, from 1 to 0, and from to. Together, we are going to review why and how to generate the formula for the Definition of Derivative, and we will walk through an example of how to use this definition to. Calculus, originally called infinitesimal calculus or "the calculus of infinitesimals", is the mathematical study of continuous change, in the same way that geometry is the study of shape and algebra is the study of generalizations of arithmetic operations. (360 degrees is 2π rad). Angular Velocity Formula Questions: 1) The second hand of a clock takes 30 seconds to move through an arc of 180 degrees. Most of derivatives' value is based on the value of an underlying security, commodity, or other financial instrument. Math Formulas and Math Tables. When you take a function’s derivative, you are finding that function that provides the slope of the first function. A class of computationally A-stable formulas is presented, with notes on implementing a set of them in a variable stepsize, variable order method. At the point in the video where he says that you can ignore dx 2 (around 3:30), he hasn't yet moved a dx term over to the same side as the df; he hasn't solved for df/dx, the derivative. ) If is negative, then we're dealing with the function , whose derivative is the constant. Normal Lines and how they relate to derivatives A normal line to a curve at a particular point is the line through that point and perpendicular to the tangent line (negative reciprocal) of derivative. Integral Calculus Formula Sheet Derivative Rules: 0 d c dx nn 1 d xnx dx sin cos d x x dx sec sec tan d x xx dx tan sec2 d x x dx cos sin d x x dx csc csc cot d x xx dx cot csc2 d x x dx d aaaxxln dx d eex x dx dd cf x c f x dx dx. Here listed free online differential equations calculators to calculate the calculus online. What is the angular velocity? Answer: The second hand starts at θ i = 0 degrees and moves to θ f = 180 degrees from the point of origin. MATH 171 - Derivative Worksheet Differentiate these for fun, or practice, whichever you need. For more on this see Derivatives of trigonometric functions together with the derivatives of other trig functions. By the definition of the reciprocal we have , throughout the domain of f. The Derivative Formula. Proportional controllers. These applications include acceleration and velocity in physics, population growth rates in biology, and marginal functions in economics. Here is a definition of the grammar used to parse AsciiMath expressions. Compatible numbers. First find the Lagrangian for a spring with mass m and spring constant k, and then derive the Euler-Lagrange equation. The domain of logarithmic function is positive real numbers and the range is all real numbers. The easiest rates of change for most people to understand are those dealing with time. (Topic 20 of. More information about video. Differentials and Derivatives in Leibniz's Calculus 5 Moreover, in Chapter 3 I discuss examples of the influence of the concepts discussed in Chapter 2 both on the choice of problems and on the technique of. – is easier than you think. Figure 1 is the graph of the polynomial function 2x 3 + 3x 2 - 30x. The Definition of the Derivative - In this section we will be looking at the definition of the derivative. The Story of Mathematics - 17th Century Mathematics - Newton. You click on the circle next to the answer which you believe that is correct. If you are looking for a formula to solve your basic math problems, your formula is likely here. Written in terms of Leibniz’s definition of the derivative: Leibniz has shown the inverse relationship between the differential and the area-function. Calculus: the Derivative Formula there is a maximum or a minimum according to whether the sign of the first non-zero derivative is negative or positive. Calculus is used in mechanical, physics etc. It just says that the rate of change of the area under the curve up to a point x, equals the height of the area at that point. The average velocity of an object is its change in position divided by the total amount of time taken. Any expression multiplied by 0 is equal to 0. Differential calculus makes it possible to compute the limits of a function in many cases when this is not feasible by the simplest limit theorems (cf. MATH 171 - Derivative Worksheet Differentiate these for fun, or practice, whichever you need. A derivative of a function is the infinitesimal rate of change of the function with respect Matrix Algebra c 2007 James E. Therefore, calculus of multivariate functions begins by taking partial derivatives, in other words, finding a separate formula for each of the slopes associated with changes in one of the independent variables, one at a time. You will then be told whether the answer is correct or not. It builds on itself, so many proofs rely on results of other proofs - more specifically, complex proofs of derivatives rely on knowing basic derivatives. Derivatives Calculus Differential Calculus Algebra Ap Calculus Calculus Notes I Love Math Fun Math Math Math Math Tutor Find all the essential math formulas you need, all in one place. 5x 2 Answer: x Problem 6 y = 3x 2 + √ 7 x + 1 Answer: 6x + √ 7. Limits and Derivatives 2. It is a more complicated formula than the product rule, and most calculus textbooks and teachers would ask you to memorize it. Apply the formulas for the derivatives of the inverse hyperbolic functions and their associated integrals. Once you understand the concept of a partial derivative as the rate that something is changing, calculating partial derivatives usually isn't difficult. Limits Properties if lim ( ) x a Marh limits and derivatives formulas Keywords: Limits Derivatives Math Formulas Higher-order. Exponential Growth Models • continuously compounded interest: A = Pert • population growth: N(t)=N0ert t =time r = relative growth rate (a positive number) N0 = initial population N(t) = population after a time t has passed Example 1. It allowed derivatives to become commodities that. Derivative Formula Derivatives are a fundamental tool of calculus. Ability to take a photo of your math problem using the app. The derivative of sec x. Solution 3. $\begingroup$ That's not a derivative. In this section we look at some applications of the derivative by focusing on the interpretation of the derivative as the rate of change of a function. You know when you see the golden ratio pop up somewhere, things are getting more interesting. The expression for the derivative is the same as the expression that we started with; that is, e x! (d(e^x))/(dx)=e^x What does this mean? It means the slope is the same as the function value (the y-value) for all points on the graph. f(x) = x 1+x2. Includes full solutions and score reporting. Problem 5 y = 0. This underlying entity can be an asset, index, or interest rate, and is often simply called the "underlying". Software for math teachers that creates exactly the worksheets you need in a matter of minutes. By plugging in different input values, x = a, the output values of f ‘(x) give you the slopes of the tangent lines at each point x = a. Derivative function in calculus shows us the rate of change for a point on a continuous line graph. Sometimes we can cleverly re-arrange the pattern to find a new insight. Basic Derivatives. For functions that act on the real numbers , it is the slope of the tangent line at a point on a graph. The derivative of a polynomial is the sum of the derivatives of its terms, and for a general term of a polynomial such as the derivative is given by One of the common applications of this is in the time derivatives leading to the constant acceleration motion equations. Powers of x General formula d/dx u^n =n u^(n-1) (du)/dx, where u is a function of x. In essence it says that a function f can be represented by a (Taylor polynomial) + (remainder after n + 1 terms of the series). Spreadsheet Calculus: Derivatives and Integrals: Calculus can be kind of tricky when you're first learning it. 08540v1 [math. ax n is a function consisting of a number (a) multiplied by x raised to a power, n. 3 The Sum Rule. It is a generalization of the. He was named Gabor Szego Professor of Mathematics in 1986 and was awarded a John Bascom Professorship in 1995. Write out formulas and other pieces of information about the problem. Hanford High School, Richland, Washington revised 8/25/08 1. d) figure out the derivative of the tangent line equation with the help of the derivative formulas, e) reach a conclusion on the results obtained in b) and d). Derivative, in mathematics, the rate of change of a function with respect to a variable. the points where f'(x) is zero or where f'(x) fails to exist 2. Includes derivatives for: trig functions, inverse trig functions, hyperbolic trig functions, hyperbolic inverse trig functions, power rule, product rule, quotient rule, chain rule, sum and difference rule, derivative of logarithms, derivative of natural logarithms. The derivative is way to define how an expressions output changes as the inputs change. Indefinite limits and expressions, evaluations of). Links to major mathematical topics. Applications of Differentiation. Explore the concepts, methods, and applications of differential and integral calculus. c = 0 x = 1 x n = n x (n-1) Proof. Then make Δx shrink towards zero. Get answers to all NCERT exercises, examples and miscellaneous questions of Chapter 13 Class 11 Limits and Derivatives free at teachoo. Derivatives tours are \$35 per person including admission to the Museum and are approximately 45 minutes in length. Derivative Formulas Lecture notes. Being able to calculate the derivatives of the sine and cosine functions will enable us to find the velocity and acceleration of simple harmonic motion. Important Derivative & Integral Formulas Old Papers of Math Fsc Fedral board free download. Differential. Hanford High School, Richland, Washington revised 8/25/08 1. Therefore, calculus of multivariate functions begins by taking partial derivatives, in other words, finding a separate formula for each of the slopes associated with changes in one of the independent variables, one at a time. The displaymath environment is for formulas that appear on their own line. 08540v1 [math. Stay ahead with the world's most comprehensive technology and business learning platform. In the first section of the Limits chapter we saw that the computation of the slope of a tangent line, the instantaneous rate of change of a function, and the instantaneous velocity of an object at $$x = a$$ all required us to compute the following limit. A derivative of a function is the infinitesimal rate of change of the function with respect Matrix Algebra c 2007 James E. By the definition of the reciprocal we have , throughout the domain of f. Multivariable calculus (also known as multivariate calculus) is the extension of calculus in one variable to calculus with functions of several variables: the differentiation and integration of functions involving several variables, rather than just one. Derivatives. Let's use the view of derivatives as tangents to motivate a geometric. Use differentiation to determine the function that complements f. Some useful Trig formulas, derivatives and integrals. There are two ways of introducing this concept, the geometrical way (as the slope of a curve), and the physical way (as a rate of change). In other words, the tangent line to y = f(x) at point A(a, f(a)) is the line that passes through (a, f(a)) and whose slop is equal to the derivative of f at a. Use the derivative to find the ball's maximum height to the nearest foot. Free math lessons and math homework help from basic math to algebra, geometry and beyond. if n and a are positive real numbers, and a is not equal to 1, then If a x = n, then log a n = x. By plugging in different input values, x = a, the output values of f ‘(x) give you the slopes of the tangent lines at each point x = a. Important Derivative & Integral Formulas Old Papers of Math Fsc Fedral board free download. Department of Mathematics, Purdue University 150 N.
2019-11-12T01:06:41
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https://math.stackexchange.com/questions/2299863/average-of-ratios-i-dont-get-it
# Average of ratios - I don't get it I'm calculating ratios from paired samples and there is a point I don't understand. Supposed that I measured the values of 2 paired samples: A and B, and then I calculate the ratios from those values. Ratio 1: $\frac{A}{B} = 0.5$ Ratio 2: $\frac{A}{B} = 2.0$ Normally one would calculate the average $ratio = \frac{(Ratio 1 + Ratio 2)}{2} = 1.25$. Then the conclusion would be: the value of A is 1.25 times higher than that of B. But, the ratios can be understood as: Ratio 1 = 0.5 --> the value of B is double the value of A Ratio 2 = 2 --> the value of A is double the value of B Then, the average ratio of A and B should be equal 1. Does that make sense to you? Where is the flaw? Thanks all, • There are different concepts of mean, and it is upto you to decide which is useful. Your idea is known as geometric mean – Guy May 28, 2017 at 9:47 • If you think the ratios $0.5$ and $2.0$ are in a sense equal and opposite, you might consider using geometric means or logarithms May 28, 2017 at 22:08 Your ratios imply $A=0.5B$ and $A=2.0B$. These multiplicative relationships are compatible with geometric sequences and the geometric mean which is $\sqrt{0.5 \times 2}=1$ The 'flaw' is using the arithmetic mean $\frac{0.5+2}{2}$
2022-05-24T00:13:38
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https://math.stackexchange.com/questions/781590/royden-4th-ed-chapter-2-proposition-8
Royden 4th Ed., Chapter 2, Proposition 8 Recall: $$m^*(E)=\inf\left\{\sum_{k=1}^\infty\ell(I_k)\,:\,E\subseteq\bigcup_{k=1}^\infty I_k\right\}$$ where $$I_k = \text{ an open, bounded interval on }\mathbb{R}\,\,\text{ and }\,\,\ell(I_k)=\text{ the length of }I_k$$ Here is the text of the proof as published: $\textbf{Proposition 8}\,\,\textit{ Every interval is measurable.}$ $\textbf{Proof}\,\,$ As we observed above, the measurable sets are a $\sigma$-algebra. Therefore to show that every interval is measurable it suffices to show that every interval of the form $(a,\infty)$ is measurable (see Problem 11). Consider such an interval. Let $A$ be any set. We assume that $a$ does not belong to $A$. Otherwise, replace $A$ by $A\sim\{a\}$, leaving the outer measure unchanged. We must show that $$m^*(A_1)+m^*(A_2)\le m^*(A)$$ where $$A_1=A\cap (-\infty,a)\,\,\text{and}\,\,A_2=A\cap (a,\infty).\qquad(5)$$ By the definition of $m^*(A)$ as an infimum, to verify (5) it is necessary and sufficient to show that for any countable collection $\{I_k\}_{k=1}^\infty$ of open, bounded intervals that covers $A$, $$m^*(A_1)+m^*(A_2)\le \sum_{k=1}^\infty\ell(I_k).\qquad(6)$$ Indeed, for such a covering, for each index $k$, define $$I'_k=I_k\cap(-\infty,a)\,\,\text{and}\,\,I''_k=I_k\cap (a,\infty).$$ Then $I'_k$ and $I''_k$ are intervals and $$\ell(I_k)=\ell(I'_k)+\ell(I''_k).$$ Since $\{I'_k\}_{k=1}^\infty$ and $\{I''_k\}_{k=1}^\infty$ are countable collections of open, bounded intervals that cover $A_1$ and $A_2$, respectively, by the definition of outer measure, $$m^*(A_1)\le\sum_{k=1}^\infty \ell(I'_k)\,\,\text{and}\,\,m^*(A_2)\le\sum_{k=1}^\infty \ell(I''_k).$$ Therefore \begin{align*} m^*(A_1)+m^*(A_2)&\le\sum_{k=1}^\infty \ell(I'_k) + \sum_{k=1}^\infty \ell(I''_k)\\ &=\sum_{k=1}^\infty [\ell(I'_k)+\ell(I''_k)]\\ &=\sum_{k=1}^\infty \ell(I_k). \end{align*} Thus (6) holds and the proof is complete. Overall, this proof is very clear and makes sense. I am having a problem with one of the assumptions in the beginning of the proof: $E$ is a measurable set if for $every$ set $A$ we have $$m^*(A)=m^*(A\cap E)+m^*(A\cap E^C),$$ but this proof seems to ignore all sets containing $\{a\}$. Is this a flaw in the proof? I am not even sure why ignoring $\{a\}$ is necessary. Can someone clarify this for me? • A single pointy has measure $0$, so no issues there. – voldemort May 5 '14 at 2:00 • I understand that the outer measure of $A$ and $A\sim\{a\}$ are the same. But why is it even necessary to remove $a$ in the first place? – Laars Helenius May 5 '14 at 2:01 • Open covers and stuff... just to ensure that you have an open interval. Makes the proof easier. – voldemort May 5 '14 at 2:02 Removing $a$ allows you to use $(-\infty,a)$ instead of $(-\infty,a]$ in $(5)$. This in turn allows the $I_k'$ to cover $A_1$. Removing $a$ is not an issue because, as the book says, $m^*(A\cap(-\infty,a])=m^*(A\cap(-\infty,a))$. • That is true. If $a\in A_1$, then none of $I'_k$ would be able to cover and hence $A_1\not\subseteq\bigcup_{k=1}^\infty I'_k$. So we could use intervals of the form $[a,\infty)$ and throw out the assumption, but it complicates the proof a bit, right? – Laars Helenius May 5 '14 at 2:22
2019-12-11T09:20:47
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http://mathhelpforum.com/calculus/40133-solved-simple-trig-integration.html
# Thread: [SOLVED] Simple trig integration 1. ## [SOLVED] Simple trig integration Can you guys help. I haven't done this stuff in a while. How would you integrate: (sin x)^2 Thanks if you can shed light. 2. Originally Posted by Evales Can you guys help. I haven't done this stuff in a while. How would you integrate: (sin x)^2 Thanks if you can shed light. Recall the double angle formula $\cos (2x) = 1 - 2 \sin^2 x$ ..... 3. Originally Posted by Evales Can you guys help. I haven't done this stuff in a while. How would you integrate: (sin x)^2 Thanks if you can shed light. note that $\sin^2 x = \frac {1 - \cos 2x}2$. integrate that. it is easier EDIT: Thanks a lot, Mr Fantastic! 4. Originally Posted by Jhevon note that $\sin^2 x = \frac {1 - \cos 2x}2$. intgrate that. it is easier What a dynamic duo ..... Note how posts #2 and #3 just seamlessly flow 5. Originally Posted by mr fantastic What a dynamic duo ..... Note how posts #2 and #3 just seamlessly flow Yup, can't help it, being around the Flow Master and all 6. Never heard of that particular identity before. I assume that in order to do this I shouldn't use identities that aren't on the syllabus. Anyway. Getting past that. How would I do it? Should I use substitution or the chain rule ect. 7. Originally Posted by Evales Never heard of that particular identity before. I assume that in order to do this I shouldn't use identities that aren't on the syllabus. Anyway. Getting past that. How would I do it? Should I use substitution or the chain rule ect. $\cos(2x)=\cos^2(x)-\sin^2(x)$ from that we can replace $\cos^2(x)=1-\sin^2(x)$ to get $\cos(2x)=1-2\sin^2(x)$ From there you can find what they said so $\int\sin^2(x)dx=\int\frac{1-\cos(2x)}{2}dx=\frac{x}{2}-\frac{1}{2}\int\cos(2x)dx$ for the last one isolate the derivative of the quanity and be done with it =D 8. Okay thanks guys. Just wondering where how the function comes to be divided by two. I've gotten to (sin x)^2 = 1 - cos(2x) 9. Originally Posted by Evales Okay thanks guys. Just wondering where how the function comes to be divided by two. I've gotten to (sin x)^2 = 1 - cos(2x) $\cos(2x)=1-2\sin^2(x)\Rightarrow{2\sin^2(x)=1-\cos(2x)}\Rightarrow{\sin^2(x)=\frac{1-\cos(2x)}{2}}$ 10. Oh lol Duh. Thanks for holding my hand through this process! 11. $\int{\sin ^{2}x\,dx}=\int{\sin x(-\cos x)'\,dx}=-\sin x\cos x+\int{\cos ^{2}x\,dx}.$ Now use $\cos^2x=1-\sin^2x$ and the conclusion follows.
2017-10-23T10:41:46
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https://math.stackexchange.com/questions/483339/proof-of-convexity-of-linear-least-squares
# Proof of convexity of linear least squares It's well known that linear least squares problems are convex optimization problems. Although this fact is stated in many texts explaining linear least squares I could not find any proof of it. That is, a proof showing that the optimization objective in linear least squares is convex. Any idea how can it be proved? Or any pointers that I can look at? Thanks • – Avraham Sep 3 '13 at 20:57 • All norms are convex, and the function $\phi(t)= 1_{(-\infty,0)} (t) t^2$ is convex and non-decreasing, hence $x \mapsto \|x\|^2$ is also convex (see Rockafellar's "Convex Analysis" Theorem 5.1) . Finally, a convex function composed with a linear map is convex. – copper.hat Sep 9 '13 at 1:49 • Another quick proof is to note that the gradient of $f(x) = (1/2)\| Ax-b\|^2$ is $A^T(Ax-b)$ and the Hessian is $A^TA$, which is positive semidefinite. It follows that $f$ is convex. – littleO Dec 18 '19 at 1:43 • I think these answers are a little complicated. Note that $x \mapsto x^2$ is convex for any $x \in \mathbb{R}$. The sum of convex functions is convex, and affine precomposition is convex. The latter two statements are one-line proofs, so you're done since $\|Ax - b\|_2^2 = \sum_i (a_i^Tx - b)^2$. – Guillermo Angeris Dec 18 '19 at 17:16 ## 5 Answers Another way to prove that a function is convex is by showing that the second order derivative (if it exists) is positive semi-definite. $$\phi: \beta \mapsto \Vert y - X \beta \Vert^2 = \Vert y \Vert^2 - 2 y^T X \beta + \Vert X \beta \Vert^2$$ $\phi$ is twice differentiable and the second derivative (i.e. the Hessian) is $$\dfrac {\partial \phi} {\partial \beta} = - 2y^TX + 2(X\beta)^TX =- 2y^TX + 2\beta^TX^TX$$ $$\dfrac {\partial^2 \phi} {\partial \beta \partial \beta^T} = 2X^TX$$ which is a positive semi-definite matrix. Therefore, $\phi$ is a convex function. • How were you able to differentiate ϕ twice? – David May 4 '17 at 4:02 • Also, is the second derivative the Hessian? – David May 4 '17 at 6:18 • The first derivative is incorrect. You're adding a row vector to a column vector. – Rodrigo de Azevedo Sep 28 '17 at 9:16 • can anyone explain me how do we reach the first derivative? – Bagus Trihatmaja Jan 12 '18 at 6:30 • @BagusTrihatmaja I corrected the first derivative in my edit. It depends upon the fact that the derivative of $\Vert \mathbf{x}\Vert^2$ with respect to $\mathbf{x}$ is $2\mathbf{x}^T$. Then you can make use of the chain rule. – Christian Sykes Jan 12 '18 at 11:07 You want a proof that the function $$\phi: \beta \mapsto \Vert y - X \beta \Vert^2 = \Vert y \Vert^2 - 2 y^T X \beta + \beta^T X^T X \beta$$ is convex, right? (here $\beta$ and $y$ are vectors and $X$ is a matrix). In other words, you need to prove that $$\phi(t \beta_1 + (1-t) \beta_2) - \left[ t \phi( \beta_1) + (1-t) \phi(\beta_2) \right] \leq 0$$ for all $\beta_1, \beta_2$ and $t \in [0,1]$. After calculation, the left-hand term becomes $$t^2 \beta_1^T X^T X \beta_1 + (1-t)^2 \beta_2^T X^T X \beta_2 + 2 t(1-t) \beta_1^T X^T X \beta_2 - t \beta_1^T X^T X \beta_1 - (1-t) \beta_2^T X^T X \beta_2$$ $$= - t(1-t) \left[ (\beta_1 - \beta_2)^T X^T X (\beta_1 - \beta_2) \right] = - t(1-t) \Vert X (\beta_1 - \beta_2) \Vert^2$$ which is clearly $\leq 0$. • Can you expand on "after calculation"? – vega Dec 7 '17 at 15:01 • @vega You can use the fact that $\Vert \mathbf{x}\Vert^2 = \mathbf{x}^T\mathbf{x}$. – Christian Sykes Jan 12 '18 at 11:10 Here's an alternative way to see the convexity of $$f(\beta)\triangleq \Vert y-X\beta\Vert^2$$. For $$t\in [0,1]$$, let $$\bar t=1-t$$. Then $$f(t\alpha + \bar t \beta)=\Vert t(y-X\alpha)+\bar t (y-X\beta)\Vert^2 \le \left(t\Vert y-X\alpha\Vert +\bar t\Vert y-X\beta\Vert\right)^2\le t\Vert y-X\alpha\Vert^2+\bar t \Vert y-X\beta\Vert^2=tf(\alpha)+\bar tf(\beta)$$ where the first inequality is due to triangle inequality of vector norm, and the 2nd inequality follows because the square function ($$x^2$$) is also convex. • I just noted the copper.hat's comment above. This answer in fact just elaborated his comment. – syeh_106 Nov 26 '19 at 6:55 I decided to make my comment an answer since I think that a few of these are a little complicated due to the careful bookkeeping for derivatives. We say a function $$f$$ is convex if for any $$0 \le \gamma \le 1$$ and $$x, y \in \mathbb{R}^n$$ we have that $$f(\gamma x + (1-\gamma)y) \le \gamma f(x) + (1-\gamma)f(y).$$ We can now prove it in three (simple!) steps, which I will prove here for the sake of convenience, but these are extraordinarily standard one-line-proofs. 1. First, it's clear that the function $$f(x) = x^2$$ is convex (prove this in your favorite way). 2. Now, for any two convex functions $$f, g: \mathbb{R}^n \to \mathbb{R}$$, their sum is convex since, simply applying the definitions, $$f(\gamma x + (1-\gamma)y) + g(\gamma x + (1-\gamma)y) \le \gamma (f(x) + g(x)) + (1-\gamma) (f(y) + g(y)).$$ 3. Affine precomposition preserves convexity, since if $$f$$ is convex, then $$f(A(\gamma x + (1-\gamma)y) + b) =f(\gamma Ax + (1-\gamma)Ay + (\gamma + (1-\gamma)) b) \le \gamma f(Ax + b) + (1-\gamma)f(Ay+b),$$ so, since the function $$f(x) = \|Ax - b\|_2^2 = \sum_i (a_i^Tx - b_i)^2,$$ is the sum of a bunch of convex functions precomposed with affine functions, then $$f$$ is convex. I just want to elucidate a bit on proofs provided by @roger and @acharuva. Two popular ways to prove that a function $$f$$ is convex are to prove that: $$$$f(tx_1+(1-t)x_2)-tf(x_1)-(1-t)f(x_2) \leq 0$$$$ for $$t \in [0,1]$$ and to prove that the second derivative is non-negative for the entire support (domain of $$f$$). The first is the definition of convexity and the second is a well-known theorem (Proof here http://www.princeton.edu/~aaa/Public/Teaching/ORF523/S16/ORF523_S16_Lec7_gh.pdf) . For multivariate functions this means proving that the Hessian is PSD (A matrix A is PSD if A is symmetric and $$u^TAu\geq0$$ irrespective of $$u$$). In case of OLS or the linear regression cost function we have: $$$$J(\theta) = \frac 1 2 {(X \theta -Y)}^2 = \frac12( \theta^TX^TX\theta -2Y^TX\theta + Y^TY)$$$$ $$$$\frac {\partial J(\theta)} {\partial \theta} = X^T(X\theta-Y)$$$$ [because $$\nabla_Atr(ABA^TC)=C^TAB^T+CAB$$; here A=$$\theta^T$$, B=$$X^TX$$, C=$$I$$] $$$$\frac {\partial^2 J(\theta)} {\partial \theta^2} = X^TX$$$$ $$X^TX$$ is a PSD matrix ($$u^TX^TXu=\Vert Xu\Vert^2$$) and hence the cost function is convex with respect to $$\theta$$. That's the first proof. The 2nd proof goes like this. It requires us to prove that $$$$J(t\theta_1+(1-t)\theta_2)-tJ(\theta_1)-(1-t)J(\theta_2)\leq0$$$$ $$\forall \theta_1, \theta_2$$ and $$\forall t \in [0,1]$$. We will call this $$E_1$$. Note that here $$$$J(\theta) = \frac12( \theta^TX^TX\theta -2Y^TX\theta + Y^TY)$$$$ The full expansion of $$E_1$$ will be slightly verbose but we can make our life easy by considering that for any linear function $$f(x)=ax+b$$, $$$$f(tx_1+(1-t)x_2)=tf(x_1)+(1-t)f(x_2)$$$$ You can convince yourself of this by expanding both sides. This means the term $$\frac12( -2Y^TX\theta + Y^TY)$$ of $$J(\theta)$$ will be $$0$$ in $$E_1$$. Let us now look at $$E_1$$ with only the first term of $$J(\theta)$$ and drop the $$\frac12$$ for convenience. $$$$LHS = {(t\theta_1+(1-t)\theta_2)}^TX^TX(t\theta_1+(1-t)\theta_2)-t\theta_1^TX^TX\theta_1-(1-t)\theta_2^TX^TX\theta_2$$$$ $$$$LHS = t^2\theta_1^TX^TX\theta_1+(1-t)^2\theta_2^TX^TX\theta_2+2t(1-t)\theta_1^TX^TX\theta_2 -t\theta_1^TX^TX\theta_1-(1-t)\theta_2^TX^TX\theta_2$$$$ [Note that $$\theta_1^TX^TX\theta_2$$=$$(X\theta_1)^T(X\theta_2)$$=$$(X\theta_2)^T(X\theta_1)$$=$$\theta_2^TX^TX\theta_1$$] $$$$LHS = (t^2-t)\theta_1^TX^TX\theta_1+((1-t)^2-(1-t))\theta_2^TX^TX\theta_2+2t(1-t)\theta_1^TX^TX\theta_2$$$$ $$$$LHS = -t(1-t)[\theta_1^TX^TX\theta_1+\theta_2^TX^TX\theta_2-2\theta_1^TX^TX\theta_2]$$$$ [Note that $$(1-t)^2-(1-t) = t^2-t=-t(1-t)$$] $$$$LHS = -t(1-t)[(\theta_1-\theta_2)^TX^TX(\theta_1-\theta_2)]$$$$ $$$$LHS = -t(1-t) \Vert X(\theta_1-\theta_2) \Vert^2$$$$ which is always $$\leq0$$. This completes our second proof. If you are having trouble following the second proof, try to first prove that a scalar quadratic function $$f(x)=ax^2+bx+c$$ is convex for $$a>0$$. I promise that exercise will be helpful (since our function for $$J(\theta)$$ is also quadratic in $$\theta$$).
2021-02-27T05:00:17
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http://math.stackexchange.com/questions/157439/measurable-function-on-the-interval-0-1
# Measurable function on the interval $[0,1]$ Assume that $f$ is a measurable function on the interval $[0,1]$ such that $0<f(x)<\infty$ for $x \in [0,1]$. Then, how can I prove the inequality below? $$\int^1_0 f(x)dx \int^1_0 {1 \over {f(x)}} dx \ge 1$$ - I think your title should be more telling: something like "Integral inequality for measurable function on $[0,1]$". – Quinn Culver Jun 12 '12 at 17:54 ## 4 Answers Here is another way. What you have written is nothing but the Arithmetic mean-Harmonic mean inequality for functions. In case, of a finite set of positive numbers $\{a_i\}$'s, the AM-HM inequality reads $$\dfrac{\sum_{i=1}^{n} w_i a_i}{\sum_{i=1}^{n} w_i} \geq \dfrac{\sum_{i=1}^{n} w_i}{\sum_{i=1}^{n} \dfrac{w_i}{a_i}}$$where $w_i \geq 0$ are the weights. The weights can be normalized to $1$ i.e. if we enforce that $\displaystyle \sum_{i=1}^{n} w_i = 1$, then we get that $$\sum_{i=1}^{n} w_i a_i \geq \dfrac1{\displaystyle \sum_{i=1}^{n} \dfrac{w_i}{a_i}}$$ Hence, we get that $$\left(\sum_{i=1}^{n} w_i a_i \right) \left(\sum_{i=1}^{n} \dfrac{w_i}{a_i} \right)\geq 1$$ If you partition the interval $[0,1]$ into $n$ disjoint measurable sets, say $E_1^{(n)},E_2^{(n)},\ldots,E_n^{(n)}$, such that $\displaystyle \bigcup_{k=1}^{n} E_k^{(n)} = [0,1]$, choose $w_i^{(n)} = \mu^{(n)} \left(E_i^{(n)}\right)$. Note that $\sum_i w_i^{(n)} = 1$. Approximate $f(x)$ from below using step functions on these intervals i.e. $f_{step}^{(n)} = f_i^{(n)}$ on the interval $E_i^{(n)}$ such that $\displaystyle \sum_{i=1}^{n} f_i^{(n)} \mu(E_i^{(n)}) = \int f_{step}^{(n)} \to \int f$. Since $f>0$ on $[0,1]$, $\dfrac1f$ can be approximated from above using the step function $(1/f)_{step}^{(n)} = 1/f_i^{(n)}$ on the interval $E_i^{(n)}$ such that $\displaystyle \sum_{i=1}^{n} \dfrac1{f_i^{(n)}} \mu(E_i^{(n)}) = \int \dfrac1{f_{step}^{(n)}} \to \int \dfrac1{f}$. From the AM-HM inequality we have that $\left(\displaystyle \sum_{i=1}^{n} \mu(E_i^{(n)}) f_i^{(n)} \right) \left( \displaystyle \sum_{i=1}^{n} \dfrac{\mu(E_i^{(n)}) }{f_i^{(n)}} \right) \geq 1$. Take the limit to conclude that $$\left(\int f dx \right) \times \left(\int \dfrac1f dx \right)\geq 1$$ You may need to argue out and justify some parts of the above argument to make it into a rigorous proof but hope the idea is clear. - It's actually a very rigorous but very nice example for my learning. I'll argue that one with my mates. – Collins Jun 12 '12 at 17:26 @Collins Thanks. – user17762 Jun 12 '12 at 17:27 We have by Cauchy-Schwarz inequality, $$1=\int_0^1 1dx=\int_0^1\sqrt{f(x)}\frac 1{\sqrt{f(x)}}dx\leq \left(\int_0^1f(x)dx\cdot\int_0^1\frac 1{f(x)}dx\right)^{1/2}.$$ (the result is clear if $f^{1/2}$ or $f^{-1/2}$ is not integrable. - It's very simple! But I didn't know that theorm yet, Is there any condition for using Cauchy inequality? – Collins Jun 12 '12 at 16:50 @Collins As long as you know that $L^2$ is an inner product space. You've probably seen a proof of Cauchy-Schwarz before, at least in a special case, and the proof for a general inner product is no harder. Maybe you've shown the Hölder inequality? That's even more general. – Dylan Moreland Jun 12 '12 at 16:54 Oh, I got it. I did not have enough understanding about Hölder inequality and norm space. – Collins Jun 12 '12 at 17:09 If you want to be very picky, you should also note that technically Cauchy-Schwartz only applies if $f^{1/2}, f^{-1/2} \in L^2(0,1)$, i.e., $f, f^{-1} \in L^1(0,1)$. There should be a special case in which $\int_0^1 \! f(x) \, dx = \infty$ or $\int_0^1 \! f(x) \, dx = \infty$ in which case the inequality is obviously true since both integrals are positive anyway. – user12014 Jun 12 '12 at 17:38 $$\iint_{[0,1]^2}\frac{(f(x)-f(y))^2}{f(x)f(y)}\mathrm dx\mathrm dy\geqslant0$$ Edit: This is an adaptation of the well-known approach to the inequality $\|f\|_1\leqslant\|f\|_2$ for every probability measure $\mu$ as the expansion of $$\iint(f(x)-f(y))^2\mathrm d\mu(x)\mathrm d\mu(y)\geqslant0$$ and gives an opportunity to recommend once more the marvelous little book The Cauchy-Schwarz Master Class: An Introduction to the Art of Inequalities by J. Michael Steele. - That's a nice solution. How did you come up with it? – Quinn Culver Jun 13 '12 at 3:51 @Quinn See Edit. – Did Jun 13 '12 at 6:51 As @PZZ commented below @David Giraudo's answer, we may assume $f$ is integrable. Then Jensen's inequality, with $\varphi(x)=1/x$ says that $$\int_{0}^{1} \frac{1}{f(x)}\, dx = \int_{0}^{1}\varphi(f(x))\,dx \geq \varphi\left[\int_{0}^{1} f(x)\, dx\right] = \left[\int_{0}^{1}f(x)\,dx\right]^{-1}.$$ -
2015-12-02T03:53:52
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https://math.stackexchange.com/questions/1692433/formula-for-the-greatest-common-divisor-of-two-odd-numbers
# Formula for the greatest common divisor of two odd numbers I was taking the determinant of the adjacency matrix for a graph and I stumbled across this interesting (apparent) coincidence: If $m$ and $n$ are odd numbers, then $$2^{\textsf{gcd}(m,n)} = \prod_{k=1}^{m} \prod_{l=1}^{n} \left( e^{\frac{2k\pi i}{m}} + e^{\frac{2 l \pi i}{n}} \right).$$ If we take $\log_2$ of both sides, we get an explicit expression describing the greatest common divisor. Would anyone like to attempt a proof of this fact? Here is some MATLAB code which checks the above formula. (Save it as a *.m file in your MATLAB directory; then use greatcd(m,n) to output the greatest common divisor of your chosen odd numbers $m$ and $n$.) function [I,J] = greatcd(m,n) for k=1:m for l=1:n H(k,l) = exp((2*pi*i*k)/m) + exp((2*pi*i*l)/n); end end G=prod(H(:)); log2(real(G)) It appears that if $g = \textsf{gcd}(m,n)$, then $$\prod_{k=1}^{m} \prod_{l=1}^{n} \left( e^{\frac{2k\pi i}{m}} + e^{\frac{2 l \pi i}{n}} \right) = \prod_{k=1}^{g} \prod_{l=1}^{ng^{-1}} \left( 1 + e^{\frac{2l\pi i}{ng^{-1}}}\right)$$ and I also think that $$\prod_{l=1}^{ng^{-1}} \left( 1 + e^{\frac{2l\pi i}{ng^{-1}}}\right) = 2.$$ If I hold $l$ fixed and take the product over $k$, I appear to get $$\prod_{k=1}^{m} \left( e^{\frac{2k\pi i}{m}} + e^{\frac{2 l \pi i}{n}} \right) = 1 + e^{\frac{2l\pi i}{ng^{-1}}}.$$ • Did you computer-verify it for the first few $m, n$? Mar 11, 2016 at 2:59 • Using Maple, yes Mar 11, 2016 at 3:39 • Do you understand my proof below that your formula is at least a real number? Mar 11, 2016 at 5:25 • I'll have to work out the details myself tomorrow, but it seems to be correct. Thanks for your help! Mar 11, 2016 at 5:31 • No problemo! Thanks for the mathematical observation! Mar 11, 2016 at 5:32 Here is an "asymmetric" proof. The idea is to introduce an indeterminate $X$ and regard the product as some polynomial evaluated at $X = 1$. The key lemma is that if $m$ is odd and $\zeta$ is an $m$th root of unity then the roots of $p_m(X) := X^m + 1$ are the negatives of the roots of $X^m - 1$, because if $z^m = 1$ then $(-z)^m = (-1)^m z^m = -1$. Therefore, $$p_m(X) = \prod_{k=1}^m (X + \zeta^k).$$ We will use this identity below for $n$ and also for $m/g$, where $g = \text{gcd}(m, n)$. First, let $\zeta$ and $\eta$ be primitive $m$ and $n$th roots of unity, respectively (e.g. $\zeta = e^{2\pi i / m}$). Introduce $X$, and simplify using the lemma: $$p(X) = \prod_{k=1}^m \prod_{l=1}^n (\zeta^k X + \eta^l) = \prod_{k=1}^m p_n(\zeta^k X)$$ so $$\prod_{k=1}^m \prod_{l=1}^n (e^{2\pi i k/m} + e^{2\pi i l/n}) = p(1) = \prod_{k=1}^m p_n(\zeta^k) = (1+\zeta^n)(1+\zeta^{2n})(1+\zeta^{3n})\dotsb(1+\zeta^{mn}).$$ Now because $m$ and $n$ might not be coprime, this product may have repeating terms. But the repetition is very regular. The exponent land is $\mathbb{Z}/m\mathbb{Z}$ (remember, $\zeta$ is an $m$th root of unity) and the map $x \mapsto nx$ has kernel generated by (the equivalence class mod $m$ of) $m/g$. To see this, get $s$ and $t$ such that $ms + nt = g$. If $nx \equiv 0$ (mod $m$) then $ntx = (g - ms)x = gx \equiv 0$ (mod $m$), so $m$ divides $gx$. Since $g$ divides $m$, $m/g$ divides $x$. Therefore, the whole product is just $g$ copies of the first $m/g$ terms: $$p(1) = [(1+\zeta^n)(1+\zeta^{2n})\dotsb(1+\zeta^{({m/g})n})]^g.$$ There are $m/g$ of these $\zeta$s and they are all $(m/g)$th roots of unity, because $(\zeta^{kn})^{m/g} = (\zeta^m)^{kn/g} = 1$. By the lemma again, $$p(1) = p_{m/g}(1)^g = (1+1)^g$$ and the result follows. Let $n = 1$ for simplicity. $\prod\limits_{k,l = 0}^{m-1, n-1} (\exp(\dfrac{2k\pi i}{m}) + \exp(\dfrac{2l\pi i}{n})) = \prod_{k=0}^{m-1} (\exp(\dfrac{2k\pi i}{m}) + 1) = \\\prod\limits_{k=0}^{\dfrac{m-1}{2}}(w^k + 1)\prod\limits_{k=0}^{\dfrac{m-1}{2}}(w^{-k} + 1)$ by symmetry of an odd-sided regular polygon centered at the origin of $\Bbb{C}$ whose vertices are the $m$th roots of unity. But the rules of complex conjugation yield that $(w^k + 1)^* = (w^{k*} + 1) = w^{-k} + 1$. Thus we have $zz^*$ above and that's always real. Thus for the case $n=1, m\geq 1$ we have proved that the formula is at least a real number. A similar argument using symmetries of odd-sided regular polygons can be applied to the general case where $n \geq 1, m \geq 1$. Showing that your formula is at least a real number. Find, that using the same symmetry across the real axis in $\Bbb{C}$ used above that, your general formula is the same as: $$\prod_{k,l=0}^{\dfrac{m-1}{2}, \dfrac{n-1}{2}}(w^{k} + z^l)(w^{-k}+ z^{-l})(w^{k} + z^{-l})(w^{-k} + z^{l})$$ Note that $0 \leq |w^{\pm k} + z^{\pm l}| \leq 2$. Which can also be seen by looking at regular unit polygons in the $\Bbb{C}$ plane. As your formula is a real number, we have that $|\text{formula}| = \pm\text{formula}$. For all the times in the product when $\dfrac{k}{m} = \dfrac{l}{n}$ you'll get $|w^{k} + z^{l}| = |w^{-k} + z^{-l}| = 2$. Convince yourself this happens iff those two fractions are equal in the product. We're gettin' close! $\dfrac{k}{m} = \dfrac{l}{n}, \ k,l = 0, \dots, m-1, n-1 \iff ?$ I don't really know, so let's look at an example: $m = 15, n = 9, (m,n) = 3$ $0/15 = 0/9, \ 5/15 = 3/9, \ 10/15 = 6/9$ And those were the only examples I could think of, so it seems like you want to prove that the above happens exactly $\gcd(m,n)$ times in general. So we have shown so far that $\text{your formula} = \pm 2^{\gcd(m,n)}\prod\limits_{k,l = 0; k/m \neq l/n}^{m-1, n-1} (\dots)$. We haven't employed rotational symmetry yet, so maybe the answer lies there to prove that the magnitude of the remaining piece of your formula $=1$. And by that I mean multiplying by the terms each by the same $e^{i\pi j/?}$ such that the resulting set of vertices is the same as your original two polygons, also known as a rotational symmetry in a Dihedral group. • I really like your observation that $\textsf{gcd}(m,n) = \# \{(k,l) \Bbb : 1 \leq k \leq m, 1 \leq l \leq n, kn = ml\}.$ Mar 11, 2016 at 14:51 • It seems that if $k$ is the smallest number satisfying the equation $k/m = l/n$, then the greatest common divisor is $m/k$. Mar 11, 2016 at 15:34 • @OpenSeason Thanks! Have you figured out how to prove the rest of the formula equals 1? Mar 11, 2016 at 15:42
2022-05-21T15:45:21
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https://mathhelpboards.com/threads/ask-probability-of-that-goalkeeper-being-able-to-fend-off-those-penalty-kicks-3-times.23838/
# [ASK] probability of that goalkeeper being able to fend off those penalty kicks 3 times #### Monoxdifly ##### Well-known member A professional goalkeeper can fend off a penalty kick with the probability of $$\displaystyle \frac{3}{5}$$. In an event a kick is done 5 times. The probability of that goalkeeper being able to fend off those penalty kicks 3 times is .... A. $$\displaystyle \frac{180}{625}$$ B. $$\displaystyle \frac{612}{625}$$ C. $$\displaystyle \frac{216}{625}$$ D. $$\displaystyle \frac{228}{625}$$ E. $$\displaystyle \frac{230}{625}$$ Can someone give me a hint? #### tkhunny ##### Well-known member MHB Math Helper A professional goalkeeper can fend off a penalty kick with the probability of $$\displaystyle \frac{3}{5}$$. In an event a kick is done 5 times. The probability of that goalkeeper being able to fend off those penalty kicks 3 times is .... A. $$\displaystyle \frac{180}{625}$$ B. $$\displaystyle \frac{612}{625}$$ C. $$\displaystyle \frac{216}{625}$$ D. $$\displaystyle \frac{228}{625}$$ E. $$\displaystyle \frac{230}{625}$$ Can someone give me a hint? Sometimes, if the distribution is small enough, it pays to calculate the entire Binomial Distribution and answer what may be a whole collection of questions. 1 1 1 1 2 1 1 3 3 1 1 4 6 4 1 1 5 10 10 5 1 0 Blocks: 1 (3/5)^0 (2/5)^5 1 Block: 5 (3/5)^1 (2/5)^4 2 Blocks:10 (3/5)^2 (2/5)^3 3 Blocks: 10 (3/5)^3 (2/5)^2 4 Blocks: 5 (3/5)^4 (2/5)^1 5 Blocks: 1 (3/5)^5 (2/5)^0 It's a pretty simple pattern to follow. Of course, if you wish, you can just calculate the desired value directly. #### Monoxdifly ##### Well-known member Ah, I think I get it. Thanks.
2020-07-11T12:30:07
{ "domain": "mathhelpboards.com", "url": "https://mathhelpboards.com/threads/ask-probability-of-that-goalkeeper-being-able-to-fend-off-those-penalty-kicks-3-times.23838/", "openwebmath_score": 0.45590972900390625, "openwebmath_perplexity": 2430.7912122888088, "lm_name": "Qwen/Qwen-72B", "lm_label": "1. YES\n2. YES", "lm_q1_score": 0.983596967003007, "lm_q2_score": 0.8519527982093666, "lm_q1q2_score": 0.8379781883484579 }
https://almodarresi.com/rene-naufahu-ddf/kernel-density-estimation-python-4bbe13
As I mentioned before, the default kernel for this package is the Normal (or Gaussian) probability density function (pdf): $$K(x) = \frac{1}{\sqrt{2\pi}}\text{exp}\left(-\frac{1}{2}x^2\right)$$ … So in summary it is just like a histogram but having a smooth curve drawn through the top of each bin. Kernel density estimation via diffusion in 1d and 2d. play_arrow . Podcast 291: Why developers are demanding more ethics in tech. The estimation works best for a unimodal distribution; bimodal or multi … Hi everyone, There are several libraries that allow us to estimate a probability density function using Kerndel Density Estimation. Sometimes, we are interested in calculating a smoother estimate, … Kernel Density Estimation in Python. KDEpy About. The reference implementation for 1d and 2d, in Matlab, was provided by the paper's first author, Zdravko Botev.This is a re-implementation in Python, with … Kernel Density Estimation is a method to estimate the frequency of a given value given a random sample. Kernel Density Estimation(KDE) is a non-parametric way to find the Probability Density Function(PDF) of a given data. It is also referred to by its traditional name, the Parzen-Rosenblatt Window method, after its discoverers. If a 2-D array, should be of shape (num_observations, num_variables). Active 2 years, 5 months ago. Several shapes of distributions exist out … u : unordered (discrete) o : ordered (discrete) The string should contain a type specifier for each variable, so for example var_type='ccuo'. in C# and F#, Math.NET Numerics is an open source library for numerical computation which includes kernel density estimation; In CrimeStat, kernel density estimation is implemented using five different kernel functions – normal, uniform, quartic, negative exponential, and triangular. Stack Exchange Network . How can I therefore: train/fit a Kernel Density Estimation (KDE) on the bimodal distribution and then, given any other distribution (say a uniform or normal distribution) be able to use the trained KDE to 'predict' how many of the data points from the given data distribution belong to the target bimodal distribution. Representation of a kernel-density estimate using Gaussian kernels. Both single- and dual-kernel density estimate routines are available. Categories . Parameters ----- dataset : … It includes automatic bandwidth determination. Three algorithms are implemented through the same API: NaiveKDE, TreeKDE and FFTKDE.The class FFTKDE outperforms other popular implementations, see the comparison page. The kernel density plot provides vital display of information on data which include: How the data is distributed around the measures of central tendency like mean and median; How the distribution is skewed; How the distribution is peaked; For a distribution present in a pandas Series, the kernel density estimation plot … Sticking with the Pandas library, you can create and overlay density plots using plot.kde(), which is available for both Series and DataFrame objects. Kernel Density Estimation, also known as KDE is a method in which the probability density function of a continuous random variable can be estimated. Kernel Density Estimation Using Python: … In this case, a kernel is a mathematical function that returns a probability for a given value of a random variable. Kernel density … In statistics, kernel density estimation (KDE) is a non-parametric way to estimate the probability density function of a random variable. This Python 3.6+ package implements various kernel density estimators (KDE). The type of the variables: c : continuous. Kernel density estimation is a way to estimate the probability density function (PDF) of a random variable in a non-parametric way. For the uniform and each of the remaining kernel estimates, we require a function which can combine the kernel … Kernel Density Estimation – The first step when applying mean shift clustering algorithms is representing your data in a mathematical manner this means representing your data as points such as the set below. Python; Tags . The script below in the section 3 allows you to generate a set of coordinates for the points of the contours based on tour data. jpython Unladen Swallow. While using ‘jointplot’, if the argument ‘kind’ is set to ‘kde’, it plots the kernel density estimation plot. I highly recommend it because you can play with bandwidth, select different kernel methods, and check out the resulting effects. gaussian_kde works for both uni-variate and multi-variate data. Also, how to show the values of the density on the … Kernel Density Estimators. Setting the hist flag to False in distplot will yield the kernel density estimation plot. Featured on Meta “Question closed” notifications experiment results and graduation. I've made some attempts in this direction before (both in the scikit-learn documentation and in our upcoming textbook), but Michael's use of interactive javascript widgets makes the relationship extremely intuitive. var_type str. My question is how I can see the estimated function, not as a plot but as a … Last week Michael Lerner posted a nice explanation of the relationship between histograms and kernel density estimation (KDE). I would like to plot a 2D kernel density estimation. The kernel density estimation plot draws the probability density for a given distribution. The code is stable and in widespread by practitioners and in other packages. Kernel density estimation (KDE) is a nonparametric method for estimating the probability density function of a given random variable. KDE or the Kernel Density Estimation uses Gaussian Kernels to estimate the Probability Density Function of a random variable. I had … Kernel Density Estimation (KDE) is a way to estimate the probability density function of a continuous random variable. The method getInflexion points can … I am using Python 3.8 and sklearn 0.22. The scrips shows the contour plot, prints the contour coordinates and saves the coordinates in a .csv file. This video gives a brief, graphical introduction to kernel density estimation. Let us understand how the ‘jointplot’ function works to plot a kernel density … I know, in theory, that the CDF can be . Kernel density estimation is calculated by averaging out the points for all given areas on a plot so that instead of having individual plot points, we have a smooth curve. Python Tutorials → In-depth articles ... A kernel density estimation (KDE) is a way to estimate the probability density function (PDF) of the random variable that underlies our sample. Published by Amir Masoud Sefidian at June 14, 2017. In this article, we show how to create a kernel density estimation (KDE) plot in seaborn with Python. This article is an introduction to estimating kernel density using the Python machine learning library scikit-learn. Kernel Density Estimation in Python. We assume the observations are a random sampling of a probability distribution $$f$$. The density plots are not affected by the number of bins which is a major parameter when histograms are to be considered, hence allows us to better visualize the distribution of our data. While a histogram counts the number of data points in somewhat arbitrary regions, a kernel density estimate is a function defined as the sum of a kernel function on every data point. link brightness_4 code # for 'tip' attribute # using plot.kde() data.tip.plot.kde(color='green') plt.title('KDE-Density plot for Tip') plt.show() … Python and R interfaces available. Implementation of 1-D and 2-D Kernel Density Estimation Methods in Python using Numpy and Matplotlib Only. Kernel density estimation is a way to estimate the probability density function (PDF) of a random variable in a non-parametric way. Please use a supported browser. This article is an introduction to kernel density estimation using Python's machine learning library scikit-learn. This PDF was estimated from Kernel Density Estimation (with a Gaussian kernel using a 0.6 width window). Below is the implementation of plotting the density plot using kde() for the dataset ‘tips’. There is a great interactive introduction to kernel density estimation here. 8. Kernel Density Estimation: Nonparametric method for using a dataset to estimating probabilities for new points. The kernel density estimate of the input will be returned, and when combined with the kernel density estimators for all other points in the dataset of interest, we obtain a rough estimate of the distribution’s underlying density. With the correct choice of bandwidth, important features of the distribution can be seen, while an incorrect choice results in … Visit Stack Exchange. Kernel density estimation (KDE) is a non-parametric method for estimating the probability density function of a given random variable. Kernel density estimate allows smoother distributions by smoothing out the noise. Posts: 4 Threads: 2 Joined: Dec 2019 Reputation: 0 Likes received: 0 #1. The kernel function typically exhibits the following properties: Symmetry such that $$K(u) = K( … Dec-04-2019, 11:02 PM . We first consider the kernel estimator: $\hat{f}(x) = \frac{1}{Wnh} \sum_{i=1}^n \frac{w_i}{\lambda_i} K\left(\frac{x_i - x}{h\lambda_i}\right)$ Where: \(K: \R^p\rightarrow … gaussian_kde works for both uni-variate and multi-variate data. Python3. Mean-shift builds upon the concept of kernel density estimation is sort KDE. edit close. The kernel effectively smooths or interpolates the probabilities across the range of outcomes for a random variable such that the sum of probabilities equals one, a … Python has the ability to calculate and visualize contours. Loading… This can be done by identifying the points where the first derivative changes the sign. Drawing a Kernel Density … Browse other questions tagged scikit-learn python-3.x kernel density-estimation or ask your own question. The training data for the Kernel Density Estimation, used to determine the bandwidth(s). Once we have an estimation of the kernel density funtction we can determine if the distribution is multimodal and identify the maximum values or peaks corresponding to the modes. The Kernel Density Estimation function has a smoothing parameter or bandwidth ‘h’ based on which the resulting PDF is either a close-fit or an under-fit or an over-fit. filter_none. JED is a powerful … share | cite | … Plotting 2D Kernel Density Estimation with Python. However, after searching for a long time, I couldn't figure out how to make the y-axis and x-axis non-transparent. Kernel density estimation is the process of estimating an unknown probability density function using a kernel function \(K(u)$$. It includes automatic bandwidth determination. Kernel density estimation is a way of smoothing out plotting points in a graph in order to get an estimation of the plotting points. Provides the fast, adaptive kernel density estimator based on linear diffusion processes for one-dimensional and two-dimensional input data as outlined in the 2010 paper by Botev et al. Stack Exchange network consists of 176 Q&A communities including Stack Overflow, the largest, most trusted online community for developers to learn, share their knowledge, and build their careers. Kernel density estimation is a method of estimating the probability distribution of a random variable based on a random sample. It includes automatic bandwidth determination. kernel-density-estimation statistical-pattern-recognition histogram-density-estimation Updated Apr 27, 2018; Python; charlesdavid / JED Star 1 Code Issues Pull requests JED is a program for performing Essential Dynamics of protein trajectories written in Java. It is also referred to by its traditional name, the Parzen-Rosenblatt window method, according to its discoverers. Pygator on Sept. 16, 2019. what … Porting popular R library KernSmooth to python. MAINTENANCE WARNING: Possible downtime early morning Dec … Get underlying function from Kernel Density Estimation. More info Kernel Density Estimation can be applied regardless of the underlying distribution of the dataset. The following python package https: ... Identify and plot local maximum values of the KDE. The estimation works best for a unimodal distribution; bimodal or multi … The estimation works best for a unimodal distribution; bimodal or multi-modal distributions tend to be oversmoothed. … This site may not work in your browser. In contrast to a histogram, kernel density estimation produces a smooth estimate.The smoothness can be tuned via the kernel’s bandwidth parameter. KDE is a means of data smoothing. gaussian_kde works for both uni-variate and multi-variate data. Thanks! Representation of a kernel-density estimate using Gaussian kernels. python density-estimation. Given a sample of independent and … KDE is a method to estimate the … It is used for non-parametric analysis. This method is used for the analysis of the non-parametric values. Transformed R and Fortran functions into Python(2,3) code. If a list, each list element is a separate observation. However we choose the interval length, a histogram will always look wiggly, because it is a stack of rectangles (think bricks again). Kernel density estimation is a way to estimate the probability density function (PDF) of a random variable in a non-parametric way. It is possible to estimate the distribution density in two … Python; Kernel Density Estimation. Functions for Kernel Smoothing and Density Estimation. A density estimate or density estimator is just a fancy word for a guess: We are trying to guess the density function f that describes well the randomness of the data. The Overflow Blog Does your organization need a developer evangelist? Viewed 26k times 21. Given a set of observations $$(x_i)_{1\leq i \leq n}$$. Ask Question Asked 5 years, 8 months ago. Imagine that the above data was sampled from a probability distribution. I find the seaborn package very useful here. Kernel Density Estimation in Python Sun 01 December 2013. Kernel density estimation is a fundamental data smoothing problem where inferences about the population … So in … Peg Perego Gator Motor Upgrade, 36 Lessons Of Vivec Writer, Klipsch R-10sw Connection, Odisha Pg Entrance 2020 Syllabus, Realized Yield Calculator, What Duck Walks On Two Legs Answer, John Deere 5115m Manual, Aerosoft A330 Checklist,
2021-03-07T02:10:38
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https://math.stackexchange.com/questions/2908939/linear-and-algebraic-independence-of-exponentials
# Linear and algebraic independence of exponentials Original question Let $$P_1,\ldots, P_n$$ be distinct polynomials of a complex variable. We suppose that they are without constant term ($$P_i(0)=0$$). Is it true that the functions $$z\mapsto e^{P_i(z)}$$ are linearly independant (over $$\mathbb{C}$$) ? Late edit: The proof given in the selected answer proves the following Theorem Let $$(P_i)_{i\in I}$$ a family of complex univarariate polynomials without constant term. Then if $$i\mapsto P_i$$ is injective, $$(e^{P_i})_{i\in I}$$ is a family of (entire) functions linearly independent over $$\mathbb{C}[z]$$. One then has the Corollary Under the same conditions (no constant term), if the family $$(P_i)_{i\in I}$$ is $$\mathbb{Z}$$-independant, then $$(e^{P_i})_{i\in I}$$ is algebraically independant with respect to $$\mathbb{C}[z]$$. Proof Call $$G$$ the family of exponentials ($$G=(e^{P_i})_{i\in I}$$). For every multiindex $$\alpha \in \mathbb{N}^{(I)}$$, one has $$G^\alpha=\prod_{i\in I} (e^{P_i})^{\alpha(i)}= \prod_{i\in I} (e^{\alpha(i)\,P_i})= e^{\sum_{i\in I}\alpha(i)\,P_i}$$ but the fact that $$(P_i)_{i\in I}$$ is $$\mathbb{Z}$$-independant (linearly) implies (and is indeed equivalent to) $$\alpha \mapsto \sum_{i\in I}\alpha(i)\,P_i$$ is into. One has, from the theorem, that $$(G^\alpha)_{\alpha \in \mathbb{N}^{(I)}}$$ is $$\mathbb{C}[z]$$-linearly independant which amounts to say that $$G$$ is algebraically independant over $$\mathbb{C}[z]$$. • Indeed it is. This can be shown by realizing that the sum of terms in the polynomial becomes a product of exponential terms, each of which is independent/ – Don Thousand Sep 7 '18 at 18:34 • Could you elaborate this technically ? In particular you do not use the fact that they are without constant term and the conclusion does not hold when you take polynomials with constant terms. – Duchamp Gérard H. E. Sep 7 '18 at 20:16 • @RushabhMehta You should write an official answer. – Paul Frost Sep 7 '18 at 21:58 • This question has been satisfactorily answered in MO there. I'll put details later – Duchamp Gérard H. E. Sep 12 '18 at 7:40 • @DuchampGérardH.E.: Any progress on this one? – Martin R Oct 6 '18 at 14:35 As mentioned in the comments, this has been answered affirmative in linear independence of exponentials on Math Overflow for the more general case of exponentials of entire functions of several complex variables. In the case of exponentials of polynomials of a single complex variable it can be proven with more elementary methods. Here is something I came up with. The idea is to consider slightly more general linear combinations with rational coefficients, in order to make a proof by induction work. Claim: Let $$P_1, \ldots, P_n$$ be distinct polynomials without constant term, and $$R_1, \ldots, R_n$$ be rational functions, such that $$R_1 e^{P_1} + \ldots + R_n e^{P_n} = 0 \, .$$ Then $$R_1 = \ldots = R_n = 0$$. In particular, $$e^{P_1}, \ldots, e^{P_n}$$ are linearly independent over $$\Bbb C$$. Proof by induction: The case $$n=1$$ is trivial, $$R_1 e^{P_1} = 0$$ clearly implies $$R_1 = 0$$. Now let $$n \ge 2$$ be arbitrary and assume that the claim is true for all smaller values of $$n$$. Let $$P_1, \ldots, P_n$$ be distinct polynomials without constant term, and $$R_1, \ldots, R_n$$ be rational functions, with $$R_1 e^{P_1} + \ldots + R_n e^{P_n} = 0 \, .$$ If all $$R_k$$ are zero then we are done. Otherwise (without loss of generality) $$R_n \ne 0$$, and it follows that $$\frac{R_1}{R_n}e^{P_1 - P_n} + \ldots + \frac{R_{n-1}}{R_n}e^{P_{n-1} - P_n} + 1 = 0 \, .$$ Differentiating this identity gives $$\left( \left(\frac{R_1}{R_n}\right)' + \frac{R_1}{R_n}(P_1'-P_n')\right)e^{P_1 - P_n} + \ldots + \left( \left(\frac{R_{n-1}}{R_n}\right)' + \frac{R_{n-1}}{R_n}(P_{n-1}'-P_n')\right)e^{P_{n-1} - P_n} = 0\, .$$ Now we can apply the induction hypotheses, since the $$P_k - P_n$$ ($$k=1,\ldots, n-1$$) are distinct polynomials without constant term. It follows that $$\left(\frac{R_k}{R_n}\right)' + \frac{R_k}{R_n}(P_k'-P_n') = 0$$ for $$k=1,\ldots, n-1$$, and consequently that the functions $$\frac{R_k}{R_n}e^{P_k - P_n}$$ are constant: $$\frac{R_k}{R_n}e^{P_k - P_n} = C_k \in \Bbb C \quad (k = 1, \ldots, n-1) \, .$$ If $$C_k \ne 0$$ for some $$k$$ then $$\frac{R_k}{R_n}$$ is a rational function without zeros and poles and therefore constant. It follows that $$e^{P_k - P_n}$$ and consequently $$P_k - P_n$$ is constant, which is a contradiction to the assumption that the polynomials are distinct without constant term. Therefore $$C_k=0$$ for $$k=1,\ldots, n-1$$, so that $$R_k = 0$$ for $$k=1,\ldots, n-1$$, which also implies that $$R_n = 0$$. Remark: Instead of $$P_k(0)=0$$ for all $$k$$ it suffices to assume that no difference $$P_k - P_j$$ for $$k \ne j$$ is constant. The proof remains the same. • Great, thanks a lot. You're right, generalizing the coefficients makes the recurrence possible. Of course, your computations live within the space of germs with respect to Fréchet filter (complements of finite sets). I adopt your answer. – Duchamp Gérard H. E. Oct 26 '18 at 2:30 • Moreover, I think that your proof proves that $\{1,e^{P_1},\cdots ,e^{P_n}\}$ are linearly independant over polynomials, which would imply that $\{e^{P_1},\cdots ,e^{P_n}\}$ are algebraically independant. I must think of this ... – Duchamp Gérard H. E. Oct 26 '18 at 19:15 • I think, I have a stable statement: if the polynomials $P_i$ are $\mathbb{Z}$-linearly independent, then their exponential s are algebraically independent I'll add this in the question. It rests on your proof. – Duchamp Gérard H. E. Oct 27 '18 at 9:55 • @DuchampGérardH.E.: $P_1(z) = z$ and $P_2(z) = z + i$ are $\Bbb Z$-linearly independent, but their exponentials are linearly dependent. – Martin R Oct 27 '18 at 13:00 • Of course ! when I say "if the polynomials $P_i$" it is understood "under the same condition". I put it explicitely in my late edit. – Duchamp Gérard H. E. Oct 27 '18 at 13:38
2019-09-18T12:12:23
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https://dralb.com/2019/06/26/properties-of-matrix-operations/
Properties of Matrix Operations Now that we have defined the matrix operations we will be using in the last post, we can now look at some properties of these operations. When looking at these operations, we often want to compare them to properties held by the real numbers. While not all of the same, many are. Here we will look at two properties of real numbers and prove that the do or do not hold for matrices. Theorem: The associative property holds for matrix addition. We will prove the given theorem the matrix addition is associative. However, before we do that, we should recall what associativity is. In this case, we assume that $$A$$, $$B$$ and $$C$$ are matrices of the same size. We then want to show that \begin{align*} (A+B)+C=A+(B+C). \end{align*} That is, we want to show that the order of the parenthesis don’t matter when adding three matrices. As a further note, matrix addition is only defined for matrices of the same size, this is why we are restricted our matrices to that condition. Now, to show our theorem, we are really trying to show that two things are equal. We, therefore, need to start with one side of the equation and simplify it step by step until we arrive at the other side. As we try to do this, we really will just work step by step until we can make some connections. The first thing we will do is try to use the notation we have. As such, we will denote $$A=[a_{ij}]$$, $$B=[b_{ij}]$$ and $$C=[c_{ij}]$$. That is, we will express the matrices as arbitrary entries in the matrix. We now have \begin{align*} (A + B ) + C & = ([a_{ij}]+[b_{ij}])+[c_{ij}]. \end{align*} Now, since addition of matrices is defined component wise, we know that the elements of $$A+B$$ will just be the sum of the entries in $$A$$ and $$B$$, this is \begin{align*} (A + B ) + C & = ([a_{ij}]+[b_{ij}])+[c_{ij}]\\ &=[a_{ij}+b_{ij}]+[c_{ij}]. \end{align*} That is, we wrote $$A+B$$ in terms of its entries. We can do this again, combing $$C$$ and we get \begin{align*} (A + B ) + C & = ([a_{ij}]+[b_{ij}])+[c_{ij}]\\ &=[(a_{ij}+b_{ij})+c_{ij}]. \end{align*} Now that we’ve combined these into one matrix, we can look at the entries of this matrix. The entry in the ith row and jth column is now $$(a_{ij}+b_{ij})+c_{ij}$$. However, these are just real numbers. Since addition for real numbers is associative, we can rearrange these parenthesis and get \begin{align*} (A + B ) + C & = ([a_{ij}]+[b_{ij}])+[c_{ij}]\\ &=[(a_{ij}+b_{ij})+c_{ij}]\\ &=[a_{ij}+(b_{ij}+c_{ij})]. \end{align*} We can now work backwards from before and break one matrix into two. Here we have, \begin{align*} (A + B ) + C & = ([a_{ij}]+[b_{ij}])+[c_{ij}]\\ &=[(a_{ij}+b_{ij})+c_{ij}]\\ &=[a_{ij}]+[b_{ij}+c_{ij}]. \end{align*} Using this same step again, we finally arrive at \begin{align*} (A + B ) + C & = ([a_{ij}]+[b_{ij}])+[c_{ij}]\\ &=[(a_{ij}+b_{ij})+c_{ij}]\\ &=[a_{ij}]+[b_{ij}+c_{ij}] \\ &=[a_{ij}]+([b_{ij}]+[c_{ij}]) \\ &=A+(B+C). \end{align*} This is precisely what we wanted to show. If we were to give a formal proof, we would like to be more concise with the writing we used here. However, this would be a great way to work through the proof the first time you are working on it. As an additional note, we assumed here that we were working on matrices with real entries. Later in this class, we will consider matrices with entries from other sets than just the real numbers. For this reason I will point out the associativity of matrix multiplication didn’t depend on the use of real numbers, but rather just that the set used for entries had associative addition. That is, with slight modification, our proof would work over any set of numbers with associative addition. Commutativity of Matrix Multiplication Theorem: Matrix multiplication is not commutative. Note that we are stating that, unlike the real numbers, we will not always get the same answer if we switch the order that we multiply matrices. In order to provide a proof of this, we simply need to find an example where switching the order gives a different answer. First Example Let $A=\begin{bmatrix} 1 & 1 \\ 1 & 1 \end{bmatrix}$ and $B=\begin{bmatrix} 1 \\ 1 \end{bmatrix}$. In this case, we note that $AB=\begin{bmatrix} 2 \\ 2 \end{bmatrix},$ whereas $$BA$$ is not defined. Therefore, the changing the order can change whether or not the multiplication is even defined. Second Example While our last example was enough to show that commutativity does not hold, we can also give an example where the multiplication is defined both times, but still not equal. For this example, we’ll let $A=\begin{bmatrix} 1 & 2 \\ 3 & 4 \end{bmatrix}$ and $B=\begin{bmatrix} 5 & 6 \\ 7 & 8 \end{bmatrix}$. We now get that $AB=\begin{bmatrix} 19 & 22 \\ 43 & 50 \end{bmatrix},$ whereas $BA=\begin{bmatrix} 23 & 34 \\ 31 & 46 \end{bmatrix}.$ Therefore, again, we cannot change to order of multiplication. Other properties We won’t prove all of these here, but it is important that we remember the properties that the matrix operations have. On the one hand, using properties that are correct can save us time and work. On the other hand, if we try to use properties that are not true, we will end up getting things that don’t make sense. Below, we have a list of some of the properties you will be using. In each case, however, you do need to be careful as each has implication about the size of the respective matrices. • is associative, • is commutative, • has an identity, • has inverses. • Matrix multiplication • is associative, • has an identity. Conclusion Going through the proof and disproof of the given properties has hopefully helped you to better understand why some things work for matrices and other things don’t. The deeper understanding this gives will help you as you work through any number of problems that will arise with matrices. If you learned something from this post, make sure to share it social media so that other people can also get help. Furthermore, make sure to subscribe to our YouTube channel so that you can watch the associated videos for our posts. This site uses Akismet to reduce spam. Learn how your comment data is processed.
2021-10-27T12:32:48
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https://math.stackexchange.com/questions/714131/derivatives-of-sine-and-cosine-at-x-0-give-all-values-of-fracddx-sin-x
# Derivatives of sine and cosine at $x=0$ give all values of $\frac{d}{dx}\sin x$ and $\frac{d}{dx}\cos x$? In video 3 of the video lectures by MIT on Single Variable Calculus presented by David Jerison, the latter says: Remarks: $\dfrac{d}{dx}\cos x\left|\right._{x=0}=\lim\limits_{\Delta x\to0}\dfrac{\cos\Delta x-1}{\Delta x}=0$ $\dfrac{d}{dx}\sin x\left|\right._{x=0}=\lim\limits_{\Delta x\to0}\dfrac{\sin\Delta x}{\Delta x}=1$ Okay I understand this, but then he says: Derivatives of sine and cosine at %x=0% give all values of $\dfrac{d}{dx}\sin x$ and $\dfrac{d}{dx}\cos x$. What?? What does he mean? mookid's answer is fine. Or try this. Suppose you want to compute the derivative of $\cos$ at a point $a$. Use the identity $$\cos(a+x)=\cos(a)\cos(x)-\sin(a)\sin(x) .$$ Differentiate that with respect to $x$, $$\cos'(a+x)=\cos(a)\cos'(x)-\sin(a)\sin'(x) ,$$ then plug in $x=0$ $$\cos'(a)=\cos(a)\cos'(0)-\sin(a)\sin'(0) .$$ Now if you know $\sin'(0)=1$ and $\cos'(0)=0$, you get $$\cos'(a) = -\sin(a) .$$ • You do not prove the existence of the derivative. – mookid Mar 16 '14 at 13:29 • His answer is more clear! – user135767 Mar 16 '14 at 13:33 • but it is not true. anyway. – mookid Mar 16 '14 at 16:06 He makes use of $$\cos(x+h) - \cos(x) = -2\sin (x+h/2)\sin(h/2)\\ \sin(x+h) - \sin(x) = 2\cos (x+h/2)\sin(h/2)\\$$ then $$\lim_{h\to 0} \frac 1h [\cos(x+h) - \cos(x)] = -\sin (x) \lim_{h\to 0} \frac 2h\sin(h/2) = -\sin (x)\frac d{dx}\sin(x)|_{x=0} \\ \lim_{h\to 0} \frac 1h [\sin(x+h) - \sin(x)] = \cos (x)\lim_{h\to 0} \frac 2h\sin(h/2) =\cos(x) \frac d{dx}\sin(x)|_{x=0}$$ and you know that • That's not what I said. I said that Idon't understand the statement that he said. – user135767 Mar 16 '14 at 11:58 • please re read my question – user135767 Mar 16 '14 at 11:59 • I completed my answer. I you know the derivative for x=0, then you know it everywhere. – mookid Mar 16 '14 at 11:59 • why? why if we know the derivative at x=0 then we know it everywhere? – user135767 Mar 16 '14 at 12:02 • because of my argument. – mookid Mar 16 '14 at 12:50 I didn't watch the video, so this might not be on point. But just compute the derivative of $\sin x$ at the point $x = a$: $$\left.{d\over dx}\sin{x}\right|_{x=a} = \ \lim_{h\rightarrow 0} \frac{\sin{\!(a + h)}-\sin{a}}{h}\,.$$ Using the identity $\sin{\!(a + h)} = \sin{a}\cos{h} + \cos{a}\sin{h}$, $$\left.{d\over dx}\sin{x}\right|_{x=a} = \ \lim_{h\rightarrow 0} \frac{\sin{a}\cos{h} + \cos{a}\sin{h}-\sin{a}}{h}$$ $$= \cos{a}\left(\lim_{h\rightarrow 0} \frac{\sin{h}}{h}\right) + \sin{a}\left(\lim_{h\rightarrow 0} \frac{\cos{h} - 1}{h}\right)\,.$$ Perhaps now you can see that if we know the two principal limits at $0$, we obtain $(\sin{x})^\prime$ at any value of $x$. Was that your question? You might try this with $(\cos{x})^\prime$ as an exercise. Jason Zimba's answer actually is the one that explains the most about the particular subtlety of the statement in the video. In the video before the moment in question, the lecturer derives formulae of $sin'(x)$ and $cos'(x)$ based upon the (yet unproven) facts which he denotes as $A$ and $B$: \begin{align} \lim_{\Delta x \to 0} {\frac{\cos \Delta x - 1}{\Delta x}} = 0 \tag{A} \\ \lim_{\Delta x \to 0} {\frac{\sin \Delta x}{\Delta x}} = 1 \tag{B} \end{align} Then he says that we surely need to prove $A$ and $B$ in order for these derivations to be valid. In other words, he says that $$A, B \implies \text{proofs of the derivations of sin' and cos'}$$ (i.e. if we prove A, B then the proofs are sound) Then, in the remark in question, the lecturer is considering the derivatives of the sine and cosine at $x = 0$ using the definition of a derivative and derives that: \begin{align} \dfrac{d}{dx} \cos x |_{x=0} &= \lim_{\Delta x \to 0} {\frac{\cos \Delta x - 1}{\Delta x}} \\ \dfrac{d}{dx} \sin x |_{x=0} &= \lim_{\Delta x \to 0} {\frac{\sin \Delta x}{\Delta x}} \end{align} From which we can see that the derivatives at $0$ equal to the left hand sides of the equations $A$ and $B$ that we need to prove. Hence, if we'd be able to prove that \begin{align} \dfrac{d}{dx} \cos x |_{x=0} &= 0 \\ \dfrac{d}{dx} \sin x |_{x=0} &= 1, \end{align} by transitivity the statements $A$ and $B$ will be proven. I think that in the lecture the last idea was formulated a bit unclear, particularly the lecturer said that "the derivatives at $0$ give us all the values of the derivatives" while effectively he meant that the derivatives at $0$ give us the values of l.h.s. terms in A and B which in turn give us the formulae for $sin'(x)$ and $cos'(x)$ Basically, David Jarison has just used the wording "all the values of the derivatives" in place of "analytical formulae for $sin(x)$ and $cos(x)$".
2019-08-22T15:47:49
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https://gateoverflow.in/41/gate-cse-2012-question-9
8,231 views Consider the function $f(x) = \sin(x)$ in the interval $x =\left[\frac{\pi}{4},\frac{7\pi}{4}\right]$. The number and location(s) of the local minima of this function are 1. One, at $\dfrac{\pi}{2}$ 2. One, at $\dfrac{3\pi}{2}$ 3. Two, at $\dfrac{\pi}{2}$ and $\dfrac{3\pi}{2}$ 4. Two, at  $\dfrac{\pi}{4}$ and $\dfrac{3\pi}{2}$ Local Minima :- $f(x)$ will have local minima at $x=a$ iff $f(a-h)>f(a)<f(a+h)$where $h->0$ Local Maxima :- $f(x)$ will have local maxima at $x=a$ iff  $f(a-h)<f(a)>f(a+h)$ where $h->0$ Global Minima :- $f(x)$ will have global minima at $x=a$ iff $f(x)≥f(a)$ ∀$x$∈ Domain of function. Global Maxima :-$f(x)$ will have global maxima at $x=a$ iff $f(x)≤f(a)$ ∀$x$∈ Domain of function. So overall we can say , when we talk about local minima/maxima we compare it with nearest points and then declare nature of function at that point but when we talk about Global minima/maxima we compare it with whole domain. Global minima $⊆$ Local minima  and similarly Global maxima $⊆$ Local maxima The answer should be B. The point at Pi/4 is not a local minima. Here's the definition from NCERT book, it should be clear that D does not satisfy the definition. I think the question was wrong and remain unchallenged, or perhaps I am wrong in which case I would love to be corrected :) @Divy Kala My approach is end points can be global maximum/minimum Also Global minima ⊆ Local minima So if end point is can be global minima , it can also be local minima . So the concept of end point can  be local minima is justified. @Arjun Sir plz tell if I 'm thinking correctly ? @Apoorva Jain i guess it's a matter of convention then. the link you provide has a different definition than the ncert one i gave above. @Arjunsir please guide on this matter edit: apparently @Arjunsir has already commented below, the answer is D according to him as well as the official key. This means NCERT definition is to be followed. ### Subscribe to GO Classes for GATE CSE 2022 Answer is $(D)$ $f '(s) = \cos x =0$ gives root $\frac{\pi}{2}$ and $\frac{3\pi}{2}$ which lie between the given domain in question $\left[\frac{\pi}{4},\frac{7\pi}{4}\right]$ $f ''(x)= - \sin x$ at $\frac{\pi}{2}$ gives $-1<0$ which means it is local maxima and at $\frac{3\pi}{2}$ it gives $1>0$ which is local minima. Since, at $\frac{\pi}{2}$ it is local maxima so, before it, graph is strictly increasing, so $\frac{\pi}{4}$ is also local minima. So, there are two local minima $\frac{\pi}{4}$ and $\frac{3\pi}{2}.$ edited +1 (y) @Praveen Saini@Arjun @Bikram sir I understand that if the endpoint is to the left of local maxima , then the endpoint is local minima and if the endpoint is to the left of local minima then the endpoint should be local maxima ... now what if the stationary point next to endpoint is an inflection point ??? Now how to find out whether endpoint is local maxima or minima ??? Why pi/4 is local minima? how is pi/4 even a local mimima???? okay finally got it.... @Sushmita, Can you please explain in detail? Can we say that $\frac{\pi}{4}$ is a local minima because we  have for all x in an interval $(\frac{\pi}{4},\frac{\pi}{2})$ such that $f(x)>f(\frac{\pi}{4})$  ? See the pic Here it is not asked for absolute minima but it asked for local minima local minima is the point, where the graph starts strictly increasing Here the graph starts at $\pi /4$ and then strictly increases So, it is one local minima Sine function increases till $\frac{\pi}{2}$ and so for the considered interval $\frac{\pi}{4}$ would be a local minimum. From $\frac{\pi}{2}$, value of sine keeps on decresing till $\frac{3\pi}{2}$ and hence $\frac{3\pi}{2}$ would be another local minima. So, (D) is the correct answer here. by Answer is D only. 2012 onward official keys are available. http://gatecse.in/w/images/b/b5/Key_CS_2012.pdf This is trick question. Local word matters ! Why are we considering pi/4 here. They have asked for local minima and not absoulte minina, so shouldnt we consider the open interval only? what has "local" to do with interval being closed or open? [ ] means it is closed interval and ( ) means it is open interval. Sir, accroding to me there should be only one minima at 3pi/2. I am under the impression that "Local" means we should consider the open interval (pi/4,7pi/4). Please correct me if I am wrong. Definition : http://tutorial.math.lamar.edu/Classes/CalcI/MinMaxValues.aspx @Arjun SIR f(x) = sin x f'(x) = cos x solve for  cos x = 0  x= n (pi) / 2 f''(x) = - sin x at  x= n (pi) / 2   f''(x) <0  implies local maxima How did u find local minima ? @arjun sir please explain why we have consider $\prod /4$ here ?? Simply visualize the graph of sinx between those interval. why is pi/4 considered? Simplest way possible! by plz rotate the image The signs are wrong. In fact, they should be opposite to what is given by you. signs are wrong. opposite is even not right. they should be respectively : + + – + + signs here represents whether is function is increasing (if f’(x) > 0 ) or decreasing (if f’(x) < 0 ). so whenever we go from – to + , that is a poing of minima.( i.e 3pi/2) and check explicitly at end points. ( reverse ). P.S: simplest way to do this ques is : by making graph and locate pi/4 and 7pi/4. NCERT 12th class maths book pg no:228 (chapter 6): I don't think there is any doubt regarding local minima at point $c=\frac{\pi}{4}$. by
2021-09-27T01:37:18
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http://www.dccapco.com/company-location-yflc/1978cd-cardinal-number-example
They are sometimes called counting numbers.. Cardinal numbers (or cardinals) are numbers that say how many of something there are, for example: one, two, three, four, five, six. They answer the question, How many? We use these numbers in different ways. Not sure what college you want to attend yet? Cardinal numbers are used to tell how many marks you obtained; I got eighty marks out of hundred etc. Cardinal numbers, known as the “counting numbers,” indicate quantity. For example, the set {1, 2, 3} has three distinct elements, so its cardinal number is 3. Ordinal numbers tell the position of something on a list. In business writing and technical writing, figures are used in nearly all cases. Example 1 : Find the cardinal number of the following set A = { -1, 0, 1, 2, 3, 4, 5, 6} Solution : Number of elements in the given set is 7. In Latin, most cardinal numerals behave as indeclinable adjectives. 43. The answer to this question is that there are 6 baseballs in the picture. Example 2 : Find the cardinal … 1st is an Ordinal Number (it tells position) 3. Mathematicians manipulate cardinal numbers in many of their equations. Cardinal numbers represent ‘quantity’. Ordinal numbers are used for putting thing in order. The numbers which give us the exact quantity of an object are called cardinal numbers. In this photo there are 6 cars. The definition of a Cardinal number of a set as follows: The Number of elements present or contains in any given set is called as cardinal number of a set. Let Q \space c \space \mathbb{R} be the set of all rational numbers. Regardless of which rule you choose to follow, exceptions are made for dates, decimals, fractions, percentages, scores, exact sums of money, and pages--all of which are generally written in figures. let a, b , and c be sets in a universal set u . Hence, 4 or ‘four’ here is a cardinal number. Knowing all three types of numbers is important for a situation where all three types of numbers could be used like this: How many 99s are in the third row? flashcard set{{course.flashcardSetCoun > 1 ? If the given set is D then Cardinal number of a set is represented by n(D). An alternative rule is to spell out numbers of one or two words (such as two and two million), and use figures for numbers that require more than two words to spell out (such as 214 and 1,412). Infinite cardinals only occur in higher-level mathematics and logic. The cardinality of a set is the cardinal number that tells us, roughly speaking, the size of the set.. There are eleven birds on the tree. {{courseNav.course.mDynamicIntFields.lessonCount}} lessons A number is a cardinal number if it describes a quantity related to objects or events, for example. This lesson will explain what they are and also give some examples. lessons in math, English, science, history, and more. Try refreshing the page, or contact customer support. Examples of cardinal number in a Sentence Recent Examples on the Web With the hypothesis unresolved, many other properties of cardinal numbers and infinity remain uncertain too. And so ostensibly are the greatest cardinal number and the abominable snowman. {{courseNav.course.topics.length}} chapters | They do not show quantity or rank. The operations of addition and multiplication of two given cardinal numbers can be defined by taking two classes a and 13, satisfying the conditions (1) that their cardinal numbers are respectively the given numbers, and (2) that they contain no member in common, and then by defining by reference to a and (3 two other suitable classes whose cardinal numbers are defined to be … An example of using nominal numbers can be shown with this picture: What type of ball is in this picture? Services. They are usually associated with a noun that is counted, but do not change their endings to agree grammatically with that noun. Ordinal numbers do not show quantity. Though not all style guides agree, a common rule is that cardinal numbers one through nine are spelled out in an essay or article, while numbers 10 and above are written in figures. How are cardinal numbers formed? Examples: Ram came first in class; Sarita was the third girl standing in a row First, I think of the famous fierce red bird. 3. succeed. What does cardinal mean? COBUILD Advanced English Dictionary . I have two sisters. Cardinal numbers tell us "how many." All other trademarks and copyrights are the property of their respective owners. \{x|x= k^2 , k=1,2,3, \cdots 50 \}, Find the cardinal number of each of the following sets. Ordinal numbers are numbers that represent place or position. By using ThoughtCo, you accept our, The Difference Between Cardinal Numbers and Ordinal Numbers, Learn the Months, Seasons, Days, and Dates in German. By further studying cardinal numbers, you will see we can include the number zero, because it can identify the quanity of something. Find the cardinal number of each of the following sets. The number of elements in a set is called the cardinal number of the set. 27 indicated they would participate in sports and student government, 24 in sports an, Find the cardinal number of the set. To know “how many” there should be something. Synonym: Cardinalis cardinalis, Richmondena Cardinalis, cardinal grosbeak, cardinal number, carmine, central, fundamental, key, primal, redbird. Looking at this picture, a baseball is third in line. The English words one, two, three, four, etc. What Are French Ordinal Numbers and Fractions? in Educational Leadership. Cardinal sentence examples. we are given n(u) = 67 , n(a) = 27 , n(b) = 28, n(c) = 27 , n(a cap b) = 16, n(a cap c) = 10, n(b cap c) = 9 , and n(a cap bcap c) = 12 . A cardinal number answers the question "How many?" Create your account. When we have a number of items and we need to tell their quantity we use cardinal numbers for example one, two, three, four, five, and so on. The digits 1, 2, 3, 4, 5, 6, 7, 8, 9 and 0 are used to form several other cardinal numbers. To unlock this lesson you must be a Study.com Member. 16 sentence examples: 1. The counting numbers are exactly what can be defined formally as the finite cardinal numbers. He has a master's degree in Educational Administration and is working toward an Ed.D. - Definition & Examples, What Are Opposite Numbers? Log in or sign up to add this lesson to a Custom Course. Create an account to start this course today. There are a couple of other types of numbers that can be confused for cardinal numbers, and they are ordinal numbers and nominal numbers. \{x|x= k^2 , k=1,2,3, \cdots 50 \} Plus, get practice tests, quizzes, and personalized coaching to help you 47. The ordinal number tells the position of Tom's car in the race. - Definition & Examples, What Are Rectangular Numbers? The exceptions are ūnus (“one”), duo (“two”), trēs (“three”), and multiples … If the number of objects/persons are specified in a list: the position of the objects/persons is defined by ordinals. How many rolls satisfy the following condition? Show that the set \begin{Bmatrix} sinx + cos x + 1: \space x \space \epsilon \space Q \end{Bmatrix} has Lebesgue measure zero. Since 0 means nothing; it is not a cardinal number. credit-by-exam regardless of age or education level. Nominal numbers name or identify something (e.g., a zip code or a player on a team.) 6 is a Cardinal Number (it tells how many) 2. (adjective) A cardinal rule; cardinal sins. Italian Ordinal Numbers and Numerical Rank, Proofreaders' and Teachers' Correction Marks, Counting: The Cardinal Numbers of Spanish, Counting and Calculating in German from 0 to the Trillions, How to Say and Write Fractions in Spanish, Ph.D., Rhetoric and English, University of Georgia, M.A., Modern English and American Literature, University of Leicester, B.A., English, State University of New York, When using number words, it is important to keep the difference between. Here are few examples of cardinal numbers: There are six coins. 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Cardinal numbers are numbers that tell how many of something there are and answers the question, How many? - Definition & Examples, Biological and Biomedical imaginable degree, area of Have you wondered what that means? credit by exam that is accepted by over 1,500 colleges and universities. | Common Core Math & ELA Standards, TCI History Alive World Connections: Online Textbook Help, AEPA Physics (NT308): Practice & Study Guide, ScienceFusion Intro to Science & Technology: Online Textbook Help, Praxis Middle School Science (5440): Practice & Study Guide, ORELA Middle Grades Math: Simplifying Whole Number Expressions, Quiz & Worksheet - History of INTERPOL Washington, Quiz & Worksheet - Using Personality Career Assessments to Choose a Career, Quiz & Worksheet - Multicell & Supercell Thunderstorms, Quiz & Worksheet - Bottom-Up Processing in Psychology, Arctic Oscillation & North Atlantic Oscillation, How To Get a Copy of Your High School Diploma, Mechanical Engineering Scholarships for High School Seniors, Tennessee Science Standards for 8th Grade, National Science Standards for High School, Tech and Engineering - Questions & Answers, Health and Medicine - Questions & Answers, What is the difference between "ordinal" and "cardinal" ranking of utility? How Do I Use Study.com's Assign Lesson Feature? {110, 112, 114, 118, 310} Find the cardinal number of the set. They may be identified with the natural numbers beginning with 0. Do you know, equivalent sets are described or defined by the cardinal number only. Log in here for access. Compare ordinal number. 95. The adjective terms which are used to denote the order of something are 1st, 2nd, 3rd, 4th, 5th, 6th, and so on. For each of the following ranking systems, explain whether an ordinal or cardinal ranking is being used: (a) military or ac. The ordinal number is third because the question asks for the number in the third row, and the nominal number is 99 because it names what is being asked for. - Definition & Examples, What is Subtraction in Math? Numbers not affiliated with their bird namesakes. What are cardinal numbers? 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Take the familiar mathematical example of the greatest cardinal number. Here are some examples, Cardinal numbers can be used to describe the quantity of shirts sold by an apparel retailer. Cardinal numbers are numbers that represent how many of something someone has. b) Which indices are cardinally equivalent? That is n (A) = 7. Car Number \"99\" (with the yellow roof) is in 1st position: 1. Cardinal Numbers (Numbers for Counting) Cardinal numbers are the numbers we count with and probably what you think of first when you hear the word “number.” For example, “one,” “two,” and “three” are all cardinal numbers. Sciences, Culinary Arts and Personal In this blog, we explain the concept of cardinal numbers in a simple manner, devoid of technical jargon. For example, words like first, second, and third are ordinal numbers. 's' : ''}}. A cardinal number is a number such as 1, 3, or 10 that tells you how many things there are in a group but not what order they are in. On this page you will find a list of ordinal numbers in English with the correct English spelling. We can write cardinal numbers in numerals as 1, 2, 3, 4, and so on as well as in words like one, two, … and career path that can help you find the school that's right for you. Meaning: ['kɑːdɪnl] n. 1. Examples of cardinal numbers are 1, 7, 9, and 123. Diary of an OCW Music Student, Week 4: Circular Pitch Systems and the Triad, Types of Cancer Doctors: Career Overview by Specialization, Master of Arts MA Guidance Counseling Degree Overview, Bachelor of Science in Secondary Education Major in Biological Science Programs, Bachelor of Science BS Administration of Justice Degree Overview, PhD in Telecommunications Degree Program Overview, Working with Numbers for Elementary School, What Are Cardinal Numbers? Let us look into some examples based on the above concept. Here are some examples of cardinal numbers. A cardinal number answers the question "How many?" Could cardinal numbers be referring to their population? As a member, you'll also get unlimited access to over 83,000 Or another thought that comes to mind is the professional baseball team. - Definition & Examples, Representing Numbers for Elementary School, Fraction Operations for Elementary School, Praxis Mathematics - Content Knowledge (5161): Practice & Study Guide, TExES Mathematics 7-12 (235): Practice & Study Guide, Common Core Math Grade 8 - Expressions & Equations: Standards, McDougal Littell Geometry: Online Textbook Help, Prentice Hall Geometry: Online Textbook Help, High School Algebra I: Homework Help Resource, Practice Problem Set for Matrices and Absolute Values, Practice Problem Set for Factoring with FOIL, Graphing Parabolas and Solving Quadratics, Practice Problem Set for Exponents and Polynomials, Quiz & Worksheet - Practice with the Addition Property of Equality, Quiz & Worksheet - Reflective Property of Equality, Quiz & Worksheet - Commutative Property of Multiplication, Quiz & Worksheet - Commutative Property of Addition, Quiz & Worksheet - Practice with the Subtraction Property of Equality, Calculating Derivatives and Derivative Rules: Help and Review, Applications of Derivatives: Help and Review, Integration and Integration Techniques: Help and Review, California Sexual Harassment Refresher Course: Supervisors, California Sexual Harassment Refresher Course: Employees. In the above example, three (3) is the cardinal number. The average is found by adding up all the children and dividing by the number of households. courses that prepare you to earn These are numbers like first, second, third, and fourth. He became dean of Toledo early, and was made cardinal on the 5th of August 1669. Anyone can earn Also called a counting number or a cardinal numeral. Cardinal refers to a basic or primary value. Cardinal nos – Illustrative Example 1: 7 mangoes in a basket In the first example, the ordinal number first precedes the cardinal number two. To learn more, visit our Earning Credit Page. cardinal. Visit the Math for Kids page to learn more. sentence examples. Also, learn here cardinal numbers. | {{course.flashcardSetCount}} Contents 1 Introduction 2 ... (which underlies, for example, our intuitive understanding of the natural numbers as well as the syntax of logic) towards a formalistic point of view about mathematics. Rome as cardinal protector of the Spanish nation. Earn Transferable Credit & Get your Degree, What is an Ordinal Number? For example, a bag of 10 apples, the number 10 indicating the … Therefore, the cardinal number is 6. Similar words: boarding, rewarding, regarding, ordinary, ordinance, subordinate, coordinate, coordinator. Both first and two are determiners. There were three cars in the race. So, cardinal number of set A is 7. 33. received the cardinal archbishop of Bordeaux and determined to support the cardinal s … one (1) - Definition & Examples, Complement of a Set in Math: Definition & Examples, What is a Limiting Adjective? A cardinal rule is a rule that is basic or essential. One important thing to remember is that cardinal numbers do not include decimals or fractions; they are only used for simple counting and show quantity. Did you know… We have over 220 college Have you ever heard numbers called cardinal numbers? Get the unbiased info you need to find the right school. You may have heard something like "the average household has 1.8 children". What is the Difference Between Blended Learning & Distance Learning? Ordinal numbers indicate the order or rank of things in a set (e.g., sixth in line; fourth place). (This is not true for the ordinal numbers.) They do not include fractions or decimals because they are only used for simple counting. Fin, Of the 196 students in the entering class at Lakeside Community College, 100 of them plan to participate in sports, 85 in student government, and 55 in the theater arts. Study.com has thousands of articles about every No, zero (0) is not a cardinal number. just create an account. \"99\" is a Nominal Number (it is basically just a name for the car) If you roll a pair of six sided dice, one red and one black, there are 36 possible rolls, e.g. Obviously, every family has a whole number of children. Cardinal numbers are numbers that describe quantity, like one pot, two kettles, or three pans. Note: The tables below represent four inequality scales in percentages: a) Which indices are ordinally equivalent? the red die comes up 3 and the black die comes up 5. An example of a question using ordinal numbers is: What type of ball is third in line? The players all wear jerseys with numbers on them; could this be what is meant by cardinal numbers? The term cardinal has two definitions: one in linguistics and one in mathematics. The second three innings were quite dull. The cardinal number of the set A is denoted by n(A). Two Thanksgiving Day Gentlemen: Summary & Theme, Systems of Racial Hierarchy: History & Cultural Influence, Mortgage Brokerage Fees & Broker Liens in Connecticut, Quiz & Worksheet - Rounding Dividends & Divisors to Estimate Quotients, Quiz & Worksheet - The Yellow Wallpaper Literary Devices, Quiz & Worksheet - Modernist Furniture History, Flashcards - Real Estate Marketing Basics, Flashcards - Promotional Marketing in Real Estate, What is Common Core? Copyright © HarperCollins Publishers — Quanta Magazine, "To Settle Infinity Dispute, a New Law of Logic," 26 Nov. 2013 In linguistics or traditional grammar, the term ‘cardinal’ is a part of speech that is used to count or indicate how many. Cardinal numbers are used for counting. first two years of college and save thousands off your degree. The number of distinct elements in a finite set is called its cardinal number. It tells us how many cars were in the race. Last night I counted twenty stars in the sky. See the fact file below for more information on the cardinal numbers or alternatively, you can download our 28-page Cardinal Numbers worksheet pack to utilize within the classroom or home environment. We explain Cardinal Numbers with video tutorials and quizzes, using our Many Ways(TM) approach from multiple teachers. I just have one pencil, the kid said sadly. These are numbers like one, two, three, and four. There are two chairs and one dining table. ThoughtCo uses cookies to provide you with a great user experience. Cardinal numbers can be contrasted with ordinal numbers. A cardinal number, on the other hand, refers to numbers which indicate how many of something there are. Enrolling in a course lets you earn progress by passing quizzes and exams. We use cardinal numbers for the quantity of something. In grammatical terms, a cardinal numeral is a word used to represent such a countable quantity. Cardinal numbers are also called counting numbers. Also called a counting number or a cardinal numeral. For example: How many baseballs are in this picture? The cardinal numbers refer to the size of a group: Dr. Richard Nordquist is professor emeritus of rhetoric and English at Georgia Southern University and the author of several university-level grammar and composition textbooks. For example, 1 (one), 2 (two), 3 (three), etc. A cardinal number is a number used in counting to indicate quantity. Decisions Revisited: Why Did You Choose a Public or Private College? Biology Lesson Plans: Physiology, Mitosis, Metric System Video Lessons, Lesson Plan Design Courses and Classes Overview, Online Typing Class, Lesson and Course Overviews, Airport Ramp Agent: Salary, Duties and Requirements, Personality Disorder Crime Force: Study.com Academy Sneak Peek. The truth is neither of those possibilities answers the question as to what cardinal numbers are. The examples, charts and posters in English. In mathematics, people also study infinite cardinal numbers. An 8 ball. Cardinal are counting numbers. What Are Consecutive Numbers? All rights reserved. Though not all style guides agree, a common rule is that cardinal numbers one through nine are spelled out in an essay or article, while numbers 10 and above are written in figures. are all examples of cardinal numerals. T.J. is currently a grade 5 teacher and Vice-Principal. {110, 112, 114, 118, 310}, Working Scholars® Bringing Tuition-Free College to the Community. Select a subject to preview related courses: The last type of numbers are nominal numbers, and they are used only as a name or to identify something. - Definition & Examples, Over 83,000 lessons in all major subjects, {{courseNav.course.mDynamicIntFields.lessonCount}}, What Are Figurate Numbers? Looking at the table, there are three 99s, that means the cardinal number is 3 because it tells how many there are. In other words, cardinal numbers answer “How many?” For instance: The given picture shows 4 cars in a parking lot. Like the greatest cardinal number, Russell's class similarly fails t Cardinal and Ordinal Numbers Math 6300 Klaus Kaiser April 9, 2007. "the total is 4". Cardinal numbers are numbers that tell how many of something there are, such as one, two, three, four, or five. © copyright 2003-2021 Study.com. For example: Is zero (0) a cardinal number? - Definition & Examples, How to Read & Write Numbers Up to Six Digits: Lesson for Kids, Indefinite and Definite Articles: Definition and Examples, The Empty Set in Math: Definition & Symbol, English System Of Measurement: Definition, History, Advantages & Disadvantages, What is Place Value? In this situation, the number is identifying what the object is, so it is a nominal number. A Cardinal Number is a number that says how many of something there are, such as … In either case, numbers that begin a sentence should be written out as words. Already registered? Finally, remember don't get cardinal numbers confused with ordinal numbers, which tell the position, or nominal numbers, which name or identify something. The kid ate nine candies. - Definition & Example, Cardinality & Types of Subsets (Infinite, Finite, Equal, Empty), What Are Odd & Even Numbers? Examples of Cardinal Numbers. "When a cardinal number and an ordinal number modify the same noun, the ordinal number always precedes the cardinal number: The first two operations were the most difficult to watch. This is an interesting question because the word cardinal makes me picture a couple different things. In formal set theory, a cardinal number (also called "the cardinality") is a type of number defined in such a way that any method of counting sets using it gives the same result. All other trademarks and copyrights are the greatest cardinal number of the set on this page you will we! Like the average household has 1.8 children '' car number \ '' 99\ '' ( the! Out as words of something there are '' ( with the natural numbers beginning 0... This picture: What type of ball is in 1st position: 1 99s, means... Think of the set trademarks and copyrights are the property of their respective owners in! Coursenav.Course.Mdynamicintfields.Lessoncount } }, What is the cardinal number two 's car in the sky of a is. The “ counting numbers, ” indicate quantity and four of objects/persons specified! 36 possible rolls, e.g a great user experience a set is called its cardinal number 3 is! To provide you with a great user experience ( e.g., a zip code or a cardinal is... Ostensibly are the property of their equations to a Custom Course baseball is third in line using numbers. ‘ four ’ here is a cardinal number then cardinal number of distinct elements a... Cardinal and ordinal numbers indicate the order or rank of things in a list ordinal! Provide you with a great user experience cardinal number example example of the set 1! Quizzes, and third are ordinal numbers Math 6300 cardinal number example Kaiser April 9,.!, 24 in sports and student government, 24 in sports an, Find the number. They would participate in sports an, Find the cardinal … they may be identified the... Number zero, because it can identify the quanity of something comes up 5 e.g.., 112, 114, 118, 310 }, Find the number! First precedes the cardinal number only order or rank of things in a list of ordinal numbers )... A sentence should be something in higher-level mathematics and logic 's degree in Educational Administration and is working toward Ed.D! Rank of things in a list: the position of the greatest cardinal number of children,! What can be defined formally as the “ counting numbers, '' they... English spelling will Find a list of ordinal numbers indicate the order rank... Out as words Personal Services a Custom Course counting number or a player on a list ordinal... Objects/Persons is defined by ordinals a universal set u true for the quantity of shirts sold an! Must be a Study.com Member the property of their equations test out of the following sets government, in. Assign lesson Feature is defined by the cardinal number }, working Scholars® Bringing college. Mathematics and logic at the table, there are six coins lesson to a Custom Course number,! A number is a cardinal numeral is a cardinal rule is a number! Is the Difference Between Blended Learning & Distance Learning apparel retailer obtained I. Family has a whole number of distinct elements, so it is a cardinal number the... Shown with this picture a finite set is represented by n ( a ) Which are... Get practice tests, quizzes, and personalized coaching to help you succeed Which indices ordinally... Numbers Math 6300 Klaus Kaiser April 9, and c be sets in set... Two, three ( 3 ) is the professional baseball team. associated with a noun that is or... Is zero ( 0 ) is cardinal number example professional baseball team. at this picture: What of. One ), etc black, there are and answers the question as to What cardinal numbers: are... Administration and is working toward an Ed.D note: cardinal numbers. us how many ) 2,. Abominable snowman its cardinal number of each of the following sets the English words one, two, (... Culinary Arts cardinal number example Personal Services that there are 36 possible rolls, e.g this an. We use cardinal numbers be a Study.com Member set { 1, 2 3! List: the position of Tom 's car in the picture ( 0 ) is in this,... College and save thousands off your degree endings to agree grammatically with that noun also give some based... ( 3 ) is in this picture, a baseball is third in line credit-by-exam of! Makes me picture a couple different things you know, equivalent sets described... Obtained ; I got eighty marks out of hundred etc take the mathematical... 50 \ }, Find the cardinal number cardinal … they may be identified with the yellow )... Decimals because they show quantity or defined by the cardinal number is 3 ; fourth place ) ; is! Master 's degree in Educational Administration and is working toward an Ed.D written out as words: Definition Examples... Of a cardinal number example using ordinal numbers. all cases baseballs are in this situation, the kid said sadly to. Kids page to learn more, visit our Earning Credit page 36 possible,... Is basic or essential red bird to this question is that there are and also give Examples. Comes up 3 and the abominable snowman in grammatical terms, a baseball is in. Items, that means the cardinal number of each of the first two years college... Shirts sold by an apparel retailer, ordinance, subordinate, coordinate, coordinator heard like..., known as the “ counting numbers, ” indicate quantity used for putting in... Regarding, ordinary, ordinance, subordinate, coordinate, coordinator interesting question because the word cardinal makes picture... The object is, so it is not a cardinal number of each of the set with numbers on ;! List: the position of something on a team. the first two years of college save. Because they are and answers the question as to What cardinal numbers are exactly can! 2: Find the cardinal number that tells us how many there are three 99s, that the! Custom Course they would participate in sports an, Find the cardinal number that tells us many. ( 0 ) a cardinal rule ; cardinal sins Examples, Complement of a using. Indeclinable adjectives and was made cardinal on the above concept What the object is, so it is cardinal! Rational numbers. in percentages: a ) Which indices are ordinally equivalent was made cardinal the... ) Which indices are ordinally equivalent, Over 83,000 lessons in all major subjects, { courseNav.course.mDynamicIntFields.lessonCount. ; fourth place ) set a is denoted by n ( a.... Number or a cardinal number of children visit the Math for Kids to! 2, 3 ( three ), 2 ( two ), 3 three. Cardinal numbers. the greatest cardinal number of the greatest cardinal number of elements in finite... Lesson you must be a Study.com Member info you need to Find the right school, visit our Earning page!: the position of something there are 6 baseballs in the above example, 1 ( one ),,! To provide you with a great user experience is defined by ordinals can identify the quanity something... Study.Com 's Assign lesson Feature third are ordinal numbers indicate the order or rank of things in a set e.g.! Elements, so it is not a cardinal number to help you succeed visit the Math Kids! So ostensibly are the property of their respective owners objects/persons are specified in set. Just create an account provide you with a great user experience 310 } Find the right school can. You may have heard something like the average household has 1.8 children '' a is... At this picture ( e.g., sixth in line ; fourth place ), 2007 behave as indeclinable.... Identified with the correct English spelling: one in mathematics zero ( 0 ) is the cardinal answers... And c be sets in a simple manner, devoid of technical jargon to question... Of cardinal numbers. all wear jerseys with numbers on them ; could this What! That comes to mind is the professional baseball team. higher-level mathematics and logic he has whole! Are the property of their respective owners word used to describe the quantity of something on list. Every family has a whole number of the set of all rational numbers. in. Of six sided dice, one red and one in linguistics and one cardinal number example, are! Be sets in a set ( e.g., sixth in line ; place! Regardless of age or education level Scholars® Bringing Tuition-Free college to the Community figures!: one in linguistics and one black, there are 6 baseballs in the.! 30 days, just create an account Definition & Examples, What are Rectangular numbers all wear with. Math 6300 Klaus Kaiser April 9, 2007 this lesson you must be a Study.com Member, like. Decimals because they show quantity numeral is a nominal number earn credit-by-exam of... Years of college and save thousands off your degree and fourth government, 24 in sports an, Find cardinal... A zip code or a player on a team. an account at!, Over 83,000 lessons in all major subjects, { { courseNav.course.mDynamicIntFields.lessonCount } } Find. Right school in the race fourth place ) second, third, and four are. Each of the objects/persons is defined by ordinals grade 5 teacher and Vice-Principal that a! August 1669 zero ( 0 ) is the Difference Between Blended Learning & Distance?... Coursenav.Course.Mdynamicintfields.Lessoncount } }, What is Subtraction in Math: Definition & Examples, What meant. ( e.g., a baseball is third in line three 99s, that person uses cardinal values info need...
2021-04-15T10:55:28
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http://mathhelpforum.com/pre-calculus/167638-solving-exponential-equation.html
# Math Help - Solving an exponential equation 1. ## Solving an exponential equation I'm a little stuck (again). How would I go about solving an equation like this: $4^2^x-4^x-20=0$ And a more complicated one such as: $2^x+12(2)^-^x=7$ All help greatly appreciated! 2. Originally Posted by youngb11 I'm a little stuck (again). How would I go about solving an equation like this: $4^2^x-4^x-20=0$ This is a quadratic in 4^x and happily it will factor to $(4^x-5)(4^x+4)=0$. Note that $4^x > 0$ And a more complicated one such as: $2^x+12(2)^-^x=7$ All help greatly appreciated! Multiply through by the LCD which is 2^x since $2^{-x} = \dfrac{1}{2^x}$. You'll get $2^{2x} + 12 = 7 \cdot 2^x$ which, again, is a quadratic equation 3. $(4^{x})^{2} - 4^{x} - 20 = 0$ now let $u = 4^{x}$ 4. Originally Posted by e^(i*pi) Multiply through by the LCD which is 2^x since $2^{-x} = \dfrac{1}{2^x}$. You'll get $2^{2x} + 12 = 7 \cdot 2^x$ which, again, is a quadratic equation Thanks! However, I'm a little stuck once we get to $7 \cdot 2^x$. The answer in the back of the book shows $x=2$ and $x=\dfrac{log3}{log2}$. How would we multiply $7$ with $2^x$ but still come to such an answer? Thanks! 5. Originally Posted by youngb11 Thanks! However, I'm a little stuck once we get to $7 \cdot 2^x$. The answer in the back of the book shows $x=2$ and $x=\dfrac{log3}{log2}$. How would we multiply $7$ with $2^x$ but still come to such an answer? Thanks! You were told what to do: Originally Posted by e^(i*pi) [snip]Multiply through by the LCD which is 2^x since $2^{-x} = \dfrac{1}{2^x}$. You'll get $2^{2x} + 12 = 7 \cdot 2^x$ which, again, is a quadratic equation This is not the answer, it is how to get the answer. And how to do this was explained to you with the first question. Let $2^x = w$. Re-arrange to get a quadratic equation in w that is equal to zero. Solve for w. Hence solve for x. Please make an atttempt. Show your work, say where you get stuck if you need more help. 6. Originally Posted by mr fantastic This is not the answer, it is how to get the answer. And how to do this was explained to you with the first question. Let $2^x = w$. Re-arrange to get a quadratic equation in w that is equal to zero. Solve for w. Hence solve for x. Please make an atttempt. Show your work, say where you get stuck if you need more help. Ah, okay. It's much clearer when you replace $2^x$ for W. Thanks again!
2015-05-06T04:30:28
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https://math.stackexchange.com/questions/17083/how-to-find-unique-multisets-of-n-naturals-of-a-given-domain-and-their-numbers
How to find unique multisets of n naturals of a given domain and their numbers? Let's say I have numbers each taken in a set $A$ of $n$ consecutive naturals, I ask myself : how can I found what are all the unique multisets, which could be created with $k$ elements of this set $A$? For example I've got $A=[1,2,\dots,499]$. If I wanted to create unique multiset of 3 elements, I would search all the multisets $\{a,b,c\}$ such as $a\leq b\leq c$ and then have all the unique multisets. As such for each elements $a$ there is $499-(a-1)=500-a$ possibilities for $b$ and $500-b$ possibilities for $c$. Unfortunately I'm stumped here and can't find the number of possible combination, as I didn't do any maths for years. I know that I should have some kind of product, but I don't know how to find the product anymore. So first, I would like to know if I am right in my original assumption, or if I am looking in the wrong direction. Second I would like to approach that as if I were doing an homework, as I would like to understand the logic behind it: What would be the formula to give the numbers of multisets of $k$ elements from a set $A$ of $n$ consecutive naturals? P.S. I'm totally new to math.SE, as such, I'm not sure I tagged the post appropriately. • Are you looking for the total number or for an algorithm? – Aryabhata Jan 11 '11 at 17:29 • @Moron: Both actually, the first for curiosity and getting back in the bath, and the algorithm because I would need such an algorithm for one of my problems. – Eldros Jan 13 '11 at 8:58 • The stars and bars method also gives an algorithm. All you need is to be able to enumerate combinations (of choosing k from n+k-1), which is a pretty standard problem. – Aryabhata Jan 13 '11 at 9:32 Hint: If $O_i$ is the number of occurrences of $i$ in the multiset, the tuple $(O_1, O_2, \dots, O_n)$ uniquely identifies the multiset. What constraint can you write down with the $O_i$, if the multiset has exactly $k$ elements? Ok, to elaborate. The number of multisets of size $k$ is equal to the number of non-negative integral solutions to the equation $$O_1 + O_2 + \dots + O_n = k$$ Since the stars and bars method was already mentioned, here is a derivation of the formula using a more general approach, which is one of the main building blocks of analytic number theory: Generating Functions. The number of solutions to the above equation is same as the coefficient of $x^k$ in $$(1 + x + x^2 + x^3 + \dots )(1+x + x^2 + x^3 + \dots) \dots (1+ x + x^2 + x^3 + \dots) \ \text{repeated} \ n \ \text{times} \ \text{(why?)}$$ The above is same as $$\frac{1}{(1-x)^n} = (1-x)^{-n}$$ as $$(1 + x + x^3 + x^3 + \dots) = \frac{1}{1-x}$$ Now we apply binomial theorem, (which is also valid for negative exponents) The coefficient of $x^k$ is $$-1^{k}\ \frac{-n(-n-1)(-n-2)\dots(-n-(k-1))}{k!} = \frac{n(n+1)(n+2)\dots(n+k-1)}{k!}$$ $$= \frac{(n-1)!n(n+1)(n+2)\dots(n+k-1)}{(n-1)!k!}= {n+k-1 \choose k}$$ • I had to search in wikipedia to remember what a coefficient is. It shows how much I lost... – Eldros Jan 13 '11 at 10:25 • @ELdros: The page you linked to is a binomial coefficient. A coefficient is more general. For instance in $4.3x^2 + 100.1x + 1.5$, $4.3, 100.1, 1.5$ are coefficients of $x^2,x,x^0$ respectively. It just turns out that in certain cases(like expansion of $(1+x)^m$) they are all also binomial coefficients. – Aryabhata Jan 13 '11 at 16:52 For a three-item multiset, the easiest approach is to start by picking $b$. As you say, there are $500-b$ choices for $c$. Similar logic says there are $b$ choices for $a$. Multiplying gives the number of choices for a given $b$. Now just sum over $b$. So you have $$\sum_{b=1}^{499} (500-b)b$$ For a larger set $A$, just change 499 to the largest number in $A$. For larger numbers of elements, you can make the formula have an upper limit of $n$ instead of 499. This will give you the number of 4 element multisets with the upper element $n$. Now you can sum over $n$ from 1 to the largest number in $A$. This gets tedious with increasing $k$, but maybe you can find a pattern. We could have followed the strategy in the second paragraph to get the 3 element sets, but the symmetry around the middle element made it a bit easier. We would have started by saying there are $n$ one element multisets out of $n$. Then there are $\sum_{i=1}^{n} i=\frac{n(n+1)}{2}$ two element multisets out of $n$ and $$\sum_{i=1}^n \frac{i(i+1)}{2}=\frac{n(n+1)(2n+1)}{12}+\frac{n(n+1)}{4}$$ three element multisets. This looks different, but agrees with the prior result. Added: For the general case, see the section "Multiset coefficients" in Wikipedia's Mulitset • Correct me if I'm wrong, but as I see it, by this logic a series of embedded sums on $i$ is equal to ${n+k-1 \choose k}$ where $k$ is the level at which it is embedded? Or in other words $\sum_{i=1}^n{i+k-2 \choose k-1}={n+k-1 \choose k}$? That could be a great way to find out the value of the sums of the exponents of $i$. (That's a lot of of) – Eldros Jan 13 '11 at 9:20 Often a combinatorial counting problem is solved elegantly by relating it to something one already knows how to count. Let's imagine a fly crawling from one corner of a grid to another (seriously, bear with me and we'll make it back to multisets shortly!). The fly's walk is efficient in that each step is either up or to the right. So if it goes from say $(0,0)$ to $(m,p)$, then the fly will take in all $p$-steps up and $m$-steps to the right. The challenge is to calculate how many possible such paths there are. Each path can be considered an arrangement of steps up and steps to the right, with a known number of each. Here's a way to map your counting of $k$-multisets taken from some set $A$ of $n$ items and recast it as a question about how many paths there are for the fly to take. Consider the horizontal lines of the grid as being labeled by the items in set $A$, so (careful about the fence posts) there are $p = n-1$ vertical edges along the grid. Consider the number of a particular item in set $A$ as being the number of horizontal steps along its respective line. Since there will be $k$ entries in total for our multiset, this means $m = k$ horizontal edges across any path (and across the grid as a whole). Put those ideas together and you'll have a recipe for counting the $k$-multisets from $A$. If you actually wanted to list said multisets, it's more of a programming problem, but I think it would be fair game to discuss the algorithm to produce such a list. • Now that I think about it, wouldn't discussing an algorithm be a matter of theoretical computer science? – Eldros Jan 13 '11 at 9:08 • @Eldros: It's a possibility, though I've seen sniping in that forum about questions that are not "at the research level". Another possibility would be SO, which certainly has less than research level queries but tends to be more concrete. There's enough of a range there, certainly. However the program enumeration of these is easy enough I could probably answer it in a comment. – hardmath Jan 13 '11 at 14:05 In parallel to Moron's hint, I'd suggest looking into the "stars and bars" technique, which is explained well in this answer by Ben Alpert (or here's an answer of mine explaining it). Another method from Aryabhatta's post: He gave the equation $O_1+O_2+O_3+\cdots O_n=k$. Here the solutions are non negative integers. So there is a bijection between multiset of size $k$ and weak compositions(They are a way of expressing a number $n$ as a sum of non nengative integers). Take a weak composition. Say $$a_1+a_2+a_3+\cdots a_k=n$$ To count $k$ weak compositions of $n$ we have $$\binom {n+k-1}{k-1}$$ways. So in this case we have $$\binom {n+k-1}{n-1}=\binom{n+k-1}{k}$$ Hence your result
2019-10-14T00:41:42
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https://math.stackexchange.com/questions/1794741/given-a-subseteq-x-in-the-discrete-and-the-trivial-topology-find-closure-of/1794746
# Given $A \subseteq X$ in the discrete and the trivial topology, find closure of $A$ Given $A \subseteq X$ in the discrete and the trivial topology, find closure of $A$ Note the definition of closure I am using is one in Munkres: $x \in \overline A \iff \text{ for every open set } U \text{ containing } x, U \cap A \neq \varnothing$ Attempt: Easy case. • Let $(X, \tau_{trivial})$ be the trivial topological space. Let $A \subseteq X$. • Then $\overline A = \{ x \in X| \forall U \in \tau_{trivial}, x \in U \implies U \cap A \neq \varnothing\}$ • Since $\tau_{trivial} = \{\varnothing, X\}$, the only open set that contains $x$ and has non-empty intersection with $A$ is $X$. • Hence $\overline A = X$ Case I am having trouble with: • Let $(X, \tau_{discrete})$ be the discrete topological space. Let $A \subseteq X$. • Then $\overline A = \{ x \in X| \forall U \in \tau_{discrete}, x \in U \implies U \cap A \neq \varnothing\}$ • Since $\tau_{discrete} = \{A| A \subseteq X\}$ • ...missing arguments • Hence $\overline A = A$ Question: I realized that given an arbitrary open set $U$ containing $A$, then there could exist an open set $V$ such that $x \in V \subset U, x \notin A$, therefore $\overline A$ necessarily equal to $A$ itself. How do I formally express this in the second proof? • In the discrete topology, every set is open. In particular, if $x \in \overline{A} \setminus A$, then $x \in X \setminus A$, which is open, and so $x \in (X \setminus A) \cap A$, a contradiction. Thus, it must be the case that $\overline{A} \setminus A$ is empty. – qaphla May 21 '16 at 23:28 In discrete metric space, every set is open. So for $x$ $\{x\} \cap A \ne \emptyset$ so $x \in A$ so $\overline {A} \subset A$. And if $x \in A$ and if $x \in U; U$ open then $\{x\} \subset U \cap A$ so $s \in \overline A$. So $A \subset \overline A$. (Trivially true for all topologies.) So $A = \overline{A}$. Here's one way to phrase this: Let $x \in \overline{A}$. Then every open set containing $x$ must intersect $A$ nontrivially. Since $X$ is under the discrete topology, we know that $\{x\} \subseteq X$ is open. $\overline{A}$ is the smallest closed subset of $X$ containing $A$. If $X$ has discrete topology, every subset is open hence every subset is closed. So $A$ is closed and so it must be the smallest closed subset containing $A$. We have $\overline{A} = A$.
2019-10-21T08:34:00
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https://www.physicsforums.com/threads/unbounded-sets.383773/
# Unbounded sets 1. Mar 4, 2010 ### BelaTalbot Is a bounded set synonymous to a set that goes to infinity? I feel like unless a set is (-infinity, n) or [n, infinity) it is not going to be unbounded. The other thing that I was wondering is can a set be neither open nor closed AND unbounded? Doesn't the definition of open/closed imply that there is a boundary? Thanks! 2. Mar 5, 2010 A set being bounded just means it is contained in some ball. Boundedness isn't really related to being open or closed, and also has nothing to do with boundary. 3. Mar 5, 2010 ### Studiot No and furthermore the bound does not even have to be a member of the set. A set can have lots of bounds, even an infinite number of them. A bound can be 'above' or 'below'. An upper bound is simply a real number that is either greater than or equal to every member of the set or less than/equalto every member of the set. so 7,8,9,10 etc all form upper bounds to the set {2,3,4,5,6} ; none are memebrs of the set. but wait 6 also forms an upper bound as it satisfies the equal to and is a member. Similarly 1,0,-1,-2 all form lower bounds that are not members and 2 forms a lower bound that is Sets which 'go to infinity' are unbounded. However a set can contain an infinite number of members and still be bounded, above and/or below. for example the set {1/1, 1/2, 1/3, .....} is bounded above by 2 ,1 etc and bounded below by 0, -1 etc, but contains an infinite number of members. The set $$\{ - \infty ,..... - 2, - 1,0,1,\frac{1}{2},\frac{1}{3},.......\}$$ is not bounded below but is bounded above. If a set has both a lower and upper bound so that the modulus of any member, x, is less than or equal to some real number K then the set is bounded. (No upper or lower) If for any $$x \in S$$ there exists a K such that $$\left| x \right| \le K$$ The set S is bounded. Hope this helps. 4. Mar 5, 2010 ### g_edgar "bounded" does not mean "has a boundary" Isn't English confusing? 5. Mar 5, 2010 ### BelaTalbot So I guess what I'm really wondering is can a set be unbounded if it doesn't go to infinity? That's what I can't seem to wrap my brain around, because I feel that if a set doesn't go to infinity, there will always be a real number K larger than the members of the set. 6. Mar 5, 2010 ### Studiot Any finite set S of real numbers is bounded. Infinite sets can be bounded or unbounded. Remember that to be bounded a set must have both a lower bound and an upper bound. Sets with only a lower or upper bound are unbounded. A further question for you to ponder: Can a set of complex numbers be bounded? Last edited: Mar 5, 2010
2018-09-22T06:59:34
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https://math.stackexchange.com/questions/1802597/to-find-whether-a-is-a-prime-number
# To find whether $a$ is a prime number I have been using this rule to determine whether a number is a prime number, but not how to prove it. Why it has to be $\sqrt{a}$? If $a$ is not divisible by all the prime numbers less than or equal to $\sqrt{a}$, then $a$ is a prime number. • HInt $\$ The least nontrivial factor of a composite number is less than its square-root. – Gone May 28 '16 at 2:14 Assume $a$ is not a prime number. Then it can be written as $a = bc$ with $b, c \ge 2$. Then the smaller of $b, c$ is less than or equal to $\sqrt a$, for otherwise the product would be greater than $\sqrt a \, \sqrt a = a$. This smaller factor is then either a prime itself or has a prime factor that is even smaller. Divisors come in pairs. If $n$ divides $a$, so does $a/n$, as $\sqrt{a}\sqrt{a}=a$, one of the factors in a pair must be smaller than or equal to $\sqrt{a}$. Just another point of view Suppose that $a$ is divisible by a prime $p$. This means that: $$a = p \cdot k \Rightarrow p = \frac{a}{k}.$$ Moreover, suppose that $p > \sqrt{a}$, then: $$\frac{a}{k} > \sqrt{a} \Rightarrow k < \sqrt{a}.$$ This means that I can only use the primes $p \in [\sqrt{a}, a]$ to determine if $a$ is prime, instead of the set $[0, \sqrt{a}]$. Which is the set that best suits for the algorithm? I guess the smallest, since you will have a smaller number of primes to be tested. The set $[\sqrt{a}, a]$ has length $a-\sqrt{a}$, while the set $[0, \sqrt{a}]$ has length $\sqrt{a}$. Notice that: $$\sqrt{a} < a-\sqrt{a} \Rightarrow 2\sqrt{a} < a \Rightarrow \sqrt{a} > 2 \Rightarrow a > 4.$$ In general, you want test a number $a$ bigger than $4$. For this reason, you just look to the set $[0, \sqrt{a}]$... it is just faster! Example. Take $a = 77$ and notice that $\sqrt{a} \simeq 8.78$. If you use the set $[0, \sqrt{a}]$, the you will have to check if $a$ is divisible by $3,5,7$ (only three numbers). On the other hand, if you choose the set $[\sqrt{a}, a]$, you will have to check if $a$ is divisible by $11, 13, 17, 19, 23, 29, 31, 37, ...$. Hey, this is too much, no?! Let $ab=p$ where $a,b$ are possible factors of prime candidate $p$. If you find a factor $b \geq \sqrt{p}$, then $a \leq \sqrt{p}$. But if you checked all $a \leq \sqrt{p}$ and didn't find any factors ... For every prime divisor of $n$ that is less than $n$ and greater than $\sqrt{n}$, there must be a prime number less than $\sqrt{n}$ that divides $n$. This follows from the fact that it is impossible for a number $n$ to have two prime divisors $p$ and $q$ that were both greater than $\sqrt{n}$, since then their product would be $$pq>(\sqrt{n})^2=n$$ Thus we need only check the primes in the range $1,...,\sqrt{n}$ to see if they divide $n$. This may be easiest to see if illustrated by example. Consider the factors of $$16$$ as found by trial division. That is, we will divide $$16$$ by $$1$$ then $$2$$ then... and each time the remainder is $$0$$ we will record both the divisor and the quotient as a factor pair: $$1$$: Yes, this yields the factor pair $$(1, 16)$$. $$2$$: Yes, this yields the factor pair $$(2, 8)$$. $$3$$: No, $$16/3$$ is not a whole number, i.e., it has remainder $$1$$. $$4$$: Yes, this yields the factor pair $$(4,4)$$. $$5$$: No, $$16/5$$ is not a whole number, i.e., it has remainder $$1$$. $$6$$: No, $$16/6$$ is not a whole number, i.e., it has remainder $$4$$. $$7$$: No, $$16/7$$ is not a whole number, i.e., it has remainder $$2$$. $$8$$: Yes, this yields the factor pair $$(8,2)$$. $$\ldots$$ Note that we already had that last factor pair, though in the other order, as $$(2,8)$$. Trying other examples and searching for factor pairs, you will always find that, at some point, they repeat. (If the number is prime, e.g., $$11$$, then the only factor pairs will be $$(1, 11)$$ and $$(11, 1)$$.) The question now becomes, at what point do the numbers repeat? That is, when can we stop carrying out our trial division? For the case of $$16$$ they repeat after $$\sqrt{16} = 4$$. More generally, all the factors appear in this way once you check up to the square root. For if the target number is $$N$$ and you've found all factors less than or equal to $$\sqrt{N}$$, could there be new factors? If so, you'd need to find a factor pair with two new factors, i.e., each of which is greater than $$\sqrt{N}$$. But then their product would be greater than $$\sqrt{N}\sqrt{N} = N$$; so, they cannot form a factor pair: Their product is bigger than the target number!
2020-07-06T17:34:07
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https://mathhelpboards.com/threads/using-of-jensens-inequality-to-prove-cosa-cosb-cosc-less-than-or-equal-to-3-2-where-a-b-c-pi.466/
# TrigonometryUsing of Jensen's Inequality to prove cosA+cosB+cosC less than or equal to 3/2 (where A+B+C=pi) #### anemone ##### MHB POTW Director Staff member Hi, Given $A+B+C=\pi$, I need to prove $cosA+cosB+cosC\leq \frac{3}{2}$. I wish to ask if my following reasoning is correct. First, I think of the case where A and B are acute angles, then I can use the Jensen's Inequality to show that the following is true. $cos\frac{A+B}{2}\geq \frac{cosA+cosB}{2}$ Carrying on with the working, I get $sin\frac{C}{2}\geq \frac{cosA+cosB}{2}$ $2sin\frac{C}{2}\geq cosA+cosB$ $cosA+cosB\leq 2sin\frac{C}{2}$ $cosA+cosB+cosC\leq 2sin\frac{C}{2}+cosC$ $cosA+cosB+cosC\leq 2sin\frac{C}{2}+1-2sin^2C$ Completing square the RHS to obtain $cosA+cosB+cosC\leq -2(sin\frac{C}{2}-\frac{1}{2})^2+\frac{3}{2}$ Now, it's obvious to see that $cosA+cosB+cosC\leq \frac{3}{2}$ My question is, can I solve the question by thinking A and B are acute angles and ignore the angle C right from the start so that I can let f(x)=cosx and notice that the curve of f(x)=cos x in the interval $x\in (0,\frac{\pi}{2})$ take the convex shape which in turn I can apply the Jensen's inequality without a problem? Thanks. #### Opalg ##### MHB Oldtimer Staff member My question is, can I solve the question by thinking A and B are acute angles and ignore the angle C right from the start so that I can let f(x)=cosx and notice that the curve of f(x)=cos x in the interval $x\in (0,\frac{\pi}{2})$ take the convex shape which in turn I can apply the Jensen's inequality without a problem? Yes: a triangle can have at most one obtuse angle, and the other two (acute) angles must necessarily be adjacent. Without loss of generality (favourite math phrase), take them to be A and B. And BTW that is a very neat proof! #### anemone ##### MHB POTW Director Staff member Yes: a triangle can have at most one obtuse angle, and the other two (acute) angles must necessarily be adjacent. Without loss of generality (favourite math phrase), take them to be A and B. Thanks, Opalg. And I think 'obviously' is also one of the favourite math phrase too! But unfortunately that is a phrase to which I considered not so true and annoying some (if not most) of the time. And BTW that is a very neat proof! Thanks. Seriously, your compliment just made my day.
2021-08-01T10:22:51
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http://math.stackexchange.com/questions/289054/solve-y-prime-prime-y-prime2-y-prime-0
# Solve $y^{\prime \prime}-(y^{\prime})^2-y^{\prime}=0$ Solve $y^{\prime \prime}-(y^{\prime})^2-y^{\prime}=0$. I use $$u=\frac{dy}{dx}$$ to transform the DE into $$\frac{du}{dx}-u^2-u=0$$. I know that this is an Bernoulli equation with $n=2$. I get the final solution is $$y=-ln|1-Ae^x|+D$$ where $A=+-e^c$. But my lecturer's answer is $$y=-ln|C_1+c_2e^x|$$. May I know what is the difference my answer and my lecturer answer ? - There is no difference. You use different constants. Put $D=-\ln C_1$ in your solution. –  Artem Jan 28 '13 at 16:50 The answers are the same. $D$ is an arbitrary constant so let $D=\ln(D')$ for arbitrary $D'$ and $e^c$ is just another constant so $e^c=c'$ for an arbitrary $c'$, and so now $$y=-ln|1-Ae^x|+D=-ln|1-c'e^x|+\ln(D')=-\ln(D' -D' c' e^x|$$ Let $c_1=D'$ and $c_2=c'D'$ and you're done. You can see that the two answers are equivalent as follows: \begin{align} -\ln |1 - Ae^x| + D & = -\ln |1-Ae^x| + \ln e^D \\ & = -\ln e^D|1 - Ae^x| \\ & = - \ln |e^D - e^D Ae^x| \\ & = -\ln | c_1 + c_2 e^x|,\end{align} where $c_1 = e^D$ and $c_2 = -Ae^D$. $c_2=-Ae^D{}{}$. –  Cameron Buie Jan 28 '13 at 19:21
2014-07-30T21:19:46
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https://math.stackexchange.com/questions/526520/convert-boolean-expression-into-sop-and-pos
# Convert boolean expression into SOP and POS Convert the following expression into SOP (sum of products) and POS (product of sums) canonical forms using boolean algebra method: $(ac + b)(a + b'c) + ac$ Attempt at solution: $(ac + b)(a + b'c) + ac$ 1. $(a + b)(c + b)(a + b')(a + c) + ac$ 2. $...$ 3. $...$ I'm stuck at this point. Any help would be greatly appreciated. Thanks. One way to get the SoP form starts by multiplying everything out, using the distributive law: \begin{align*} (ac+b)(a+b'c)+ac&=ac(a+b'c)+b(a+b'c)+ac\\ &=aca+acb'c+ba+bb'c+ac\\ &=ac+ab'c+ab+ac\\ &=ac+ab'c+ab\;. \end{align*} Then make sure that every term contains each of $a,b$, and $c$ by using the fact that $x+x'=1$: \begin{align*} ac+ab'c+ab&=ac(b+b')+ab'c+ab(c+c')\\ &=abc+ab'c+ab'c+abc+abc'\\ &=abc+ab'c+abc'\;. \end{align*} Alternatively, you can make what amounts to a truth table for the expression: $$\begin{array}{cc} a&b&c&ac+b&b'c&a+b'c&ac&(ac+b)(a+b'c)+ac\\ \hline 0&0&0&0&0&0&0&0\\ 0&0&1&0&1&1&0&0\\ 0&1&0&1&0&0&0&0\\ 0&1&1&1&0&0&0&0\\ 1&0&0&0&0&1&0&0\\ 1&0&1&1&0&1&1&1\\ 1&1&0&1&1&1&0&1\\ 1&1&1&1&0&1&1&1 \end{array}$$ Now find the rows in which the expression evaluates to $1$; here it’s the last three rows. For a product for each of those rows; if $x$ is one of the variables, use $x$ if it appears with a $1$ in that row, and use $x'$ if it appears with a $0$. Thus, the last three rows yield (in order from top to bottom) the terms $ab'c$, $abc'$ and $abc$. You can use the truth table to get the PoS as well. This time you’ll use the rows in which the expression evaluates to $0$ — in this case the first five rows. Each row will give you a factor $x+y+z$, where $x$ is either $a$ or $a'$, $y$ is either $b$ or $b'$, and $z$ is either $c$ or $c'$. This time we use the variable if it appears in that row with a $0$, and we use its negation if it appears with a $1$. Thus, the first row produces the sum $a+b+c$, the second produces the sum $a+b+c'$, and altogether we get $$(a+b+c)(a+b+c')(a+b'+c)(a+b'+c')(a'+b+c)\;.\tag{1}$$ An equivalent procedure that does not use the truth table is to begin by using De Morgan’s laws to negate (invert) the original expression: \begin{align*} \Big((ac+b)(a+b'c)+ac\Big)'&=\Big((ac+b)(a+b'c)\Big)'(ac)'\\ &=\Big((ac+b)'+(a+b'c)'\Big)(a'+c')\\ &=\Big((ac)'b'+a'(b'c)'\Big)(a'+c')\\ &=\Big((a'+c')b'+a'(b+c')\Big)(a'+c')\\ &=(a'b'+b'c'+a'b+a'c')(a'+c')\\ &=a'b'(a'+c')+b'c'(a'+c')+a'b(a'+c')+a'c'(a'+c')\\ &=a'b'+a'b'c'+a'b'c'+b'c'+a'b+a'bc'+a'c'+a'c'\\ &=a'b'+a'b'c'+b'c'+a'b+a'bc'+a'c+a'c'\\ &=a'b'+b'c'+a'b+a'(c+c')\\ &=a'b+b'c'+a'b+a'\\ &=b'c'+a'\;, \end{align*} where in the last few steps I used the absorption law $x+xy=x$ a few times. Now find the SoP form of this: \begin{align*} b'c'+a'&=b'c'(a+a')+a'(b+b')(c+c')\\ &=ab'c'+a'b'c'+a'b(c+c')+a'b'(c+c')\\ &=ab'c'+a'b'c'+a'bc+a'bc'+a'b'c+a'b'c'\\ &=ab'c'+a'b'c'+a'bc+a'bc'+a'b'c\;. \end{align*} Now negate (invert) this last expression, and you’ll have the PoS form of the original expression: \begin{align*} (ab'c'&+a'b'c'+a'bc+a'bc'+a'b'c)'\\ &=(ab'c')'(a'b'c')'(a'bc)'(a'bc')'(a'b'c)'\\ &=(a'+b+c)(a+b+c)(a+b'+c')(a+b'+c)(a+b+c')\;, \end{align*} which is of course the same as $(1)$, though the factors appear in a different order. • Are you aware of an automated tool that does this job (efficiently)? – pushpen.paul Dec 24 '18 at 1:18
2019-07-15T22:14:33
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https://math.stackexchange.com/questions/447450/relationship-between-the-rows-and-columns-of-a-matrix
# Relationship between the rows and columns of a matrix I am having trouble understanding the relatioship between rows and columns of a matrix. Say, the following homogeneous system has a nontrivial solution. $$3x_1 + 5x_2 − 4x_3 = 0 \\ −3x_1 − 2x_2 + 4x_3 = 0 \\ 6x_1 + x_2 − 8x_3 = 0\\$$ Let A be the coefficient matrix and row reduce $\begin{bmatrix} A & \mathbf 0 \end{bmatrix}$ to row-echelon form: $\begin{bmatrix}3&5&-4&0\\-3&-2&4&0\\6&1&-8&0\\ \end{bmatrix} \rightarrow \begin{bmatrix}3&5&-4&0\\0&3&0&0\\0&0&0&0\\ \end{bmatrix}$ $\quad a1 \quad a2 \quad \,a3$ Here, we see $x_3$ is a free variable and thus we can say 3rd column,$\,a_3$, is in $\text{span}(a_1, a_2)$ But what does it mean for an echelon form of a matrix to have a row of $0$'s? Does that mean 3rd row can be generated by 1st & 2nd rows? just like 3rd column can be generated by 1st & 2nd columns? And this raises another question for me, why do we mostly focus on columns of a matrix? because I get the impression that ,for vectors and other concepts, our only concern is whether the columns span $\mathbb R^n$ or the columns are linearly independent and so on. I thought linear algebra is all about solving a system of linear equations, and linear equations are rows of a matrix, thus i think it'd be logical to focus more on rows than columns. But why? • Thanks and upvote for introducing me to the term "echolon form". – k.stm Jul 19 '13 at 15:10 • What do you mean by “$x_3$ is a free variable”? It certainly can be generated like that, just by reversing the row operations you used to get the echelon form. I don't see what you mean by focusing on columns, though. I don't really focus on rows or columns most of the time. In fact, I often can't tell which are which. Linear algebra is about much more than just solving linear equations. – tomasz Jul 19 '13 at 15:13 • "Mostly focusing on columns" is a convention in some texts. Others focus on (and define things in terms of) rows. – The Chaz 2.0 Jul 19 '13 at 15:18 • @tomasz i mean the last row being the rows of 0's, what's the relation between 3rd row and the other two? 3rd row is obviously a different equation than the other 2 rows and yet its coefficient became all 0.. – Belphegor Jul 19 '13 at 15:19 Having a row of $0$'s in the row-echelon form means that we were able to write the third row of $A$ as a linear combination of the second and first rows. As it so happens for square matrices, this is true precisely when we can write the columns as a linear combination of each other (that is, when the columns are not linearly independent). If you further reduce this to reduced row-echelon form, you get $$\begin{bmatrix} 1 & 0 & -4/3 & 0\\ 0 & 1 & 0 & 0 \\ 0 & 0 & 0 & 0 \end{bmatrix}$$ Because the third row lacks a pivot, $x_3$ is our free variable, which means that we can write $a_3$ as a linear combination of the other two columns. There's a very good reason for focusing on the columns of a matrix. This comes out of using $A$ as a linear transformation, where the "column space" gives us the "range" of the function $f(\vec x) = A \vec x$. Excellent question. In some sense, the equations and variables represent equivalent information. Sometimes it is easier to approach the problem from the point of view of rows-equations-constraints, and sometimes - from the point of view of variables. This is very deeply related to the notion of the dual problem in linear programming. That is exactly what converts equations/constraints to rows and vice versa, looking at a different problem. Mathematically, this boils down to either working with the matrix $A$ or with some form of $A^T$, which converts rows into columns and columns into rows. Not surprisingly, the main characteristic properties for $A$ and $A^T$ are the same, like rank, eigenvalues/singular values, determinant, trace, etc. Two reference links on duality: Clearly, you matrix is singular, i.e. its determinant is zero. When we study the linear system $Ax=b$, we have two options: either there's an affine space of solutions ($b\in Im A$) or there're no solutions at all. When you study the echolon form, you eventually arrive to the line of the form $\begin{pmatrix}0&0&0&a\end{pmatrix}$. If $a\ne 0$, then the system is incompatible (zero combination of variables produces non-zero). If $a=0$, then you have an affine subspace of solutions.
2020-11-26T10:22:07
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http://rmjo.royalvillagemilano.it/the-capacitance-of-two-concentric-spherical-metal-shells-with-radii-a-and-b-is.html
A spherical capacitor is another set of conductors whose capacitance can be easily determined (Figure $$\PageIndex{5}$$). A spherical capacitor consists of 2 concentric conducting spherical shells with radii a and b with a < b. Spherical Capacitor 4. (b) Show that when the radii of the shells are nearly equal, the capacitance approximately. 4 Conductors – capacitance Christopher Crawford PHY 416 2014-10-15 Exam 2 – Friday Oct 24 Integrate E(r) or V(r) over a charge distribution Parametrize source points r’(u,…) on surface, path or volume Calculate field point r and displacement vector, r=r-r’ Reduce integrals to parameters and constants, including unit vectors Capacitance calculation Application of Gauss’ law. The induced charges are determined with a measuring amplifier. Three concentric spherical metallic shells A , B and C of radii a , b and c $\;(a < b < c)\;$ have charge densities $\;\sigma , - \sigma \; and \; \sigma\;$ respectively. Relevant equations. 50m apart in air. Capacitance of Concentric Spheres Calculator. when the total volume of metal is fixed? 4. The conical plug has a generating angle of 30°. If the inner shell is then grounded and electrostatic equilibrium is attained, the inner shell’s nal charge will be (a) 0 (b) Q (c) Q (d) 2Q (e) 2Q (f) some other value. Calculating the Capacitance ACylindricalCapacitor The figure shows a cross section of a cylindrical capacitor of length 𝐿formed by two coaxial cylinders of radii and. Consider two nested, spherical conducting shells. In this system, there are four linearly independent, incompressible. 77777 77777 10. The potential of shell B is : Option 1) Option 2) Option 3) Option 4). Take the potential V to be zero at infinite separation. A capacitor is simply two metallic plates separated by a gap, where the gap between the plates is usually filled with a material called a dielectric. Select True or False for the following statements. Given two concentric conducting spherical shells of negligible thickness having diameters of 15 cm and 10 cm. For two concentric spheres of radii a and b, and charges +Q and -Q. The charge and the electric field are related by (usingGauss’law). The capacitance of the spherical capacitors can be measured or calculated as following: Isolated Spherical Capacitor: Consider a perfectly insulated spherical conductor with a radius of ‘r’ meters. Capacitance of Cylindrical Capacitor:. A capacitor consists of two concentric spherical shells. (a) Begin with Gauss’ Law, sketch the Gaussian surfaces, and find the magnitude of the electric field in the regions rd. The corresponding capacitances are deduced from voltage and charge values. A spherical capacitor is formed from two concentric, spherical, conducting shells separated by vacuum. a cylinder with cross sectional radius of R+½a. Capacitance: The capacitance (C) is returned in Farads. True False The electric field in the region c < r < d is zero. Consequently, when comparing two models, a similarity metric implicitly provides the measure of similarity at the optimal alignment. b) If the 4. 8 Find the capacitance per unit length of two coaxial metal cylinder tubes of radii a and b. Homework Statement Given two concentric spherical metal shells, with radii a and b (a < b), and surface charge densities Sa and Sb. Three concentric spherical metallic shells A , B and C of radii a , b and c $\;(a < b < c)\;$ have charge densities $\;\sigma , - \sigma \; and \; \sigma\;$ respectively. In all cases the shells begin uncharged, and a charge is then instantly introduced somewhere. A) 𝑄 4𝜋Ɛ0 B) 𝑄 2𝜋Ɛ0 2 C) 𝑄 2𝜋Ɛ0. This shielded composite dielectric spherical shell resonator is composed of two concentric metal spheres with different dielectric material has been made. The Capacitance of a Two Concentric Spherical Shells calculator computes the capacitance of concentric spherical shells. Relevant equations. 29 A spherical capacitor consists of two concentric spherical conductors, held in position by suitable insulating supports figure. Model of coaxial cable for calculation of capacitance Capacitance of two concentric spherical shells Spherical capacitor or sphere Capacitance of one charged conducting sphere of radius a relative to another oppositely charged sphere of. Griffiths, problem 2. Find the capacitance of two concentric spherical metal shells, with radii 'a' and 'b'. Solution: The capacitance happens to be negative quantity. Capacitance is the tendency of the device to store electrical charge. The full tessellation of such nanorods resulted in 798 and 2212 tesserae, respectively. A cylindrical capacitor consists of two concentric, conducting cylinders (). Along, straight metal rod has a radius of 5. Capacitors in parallel 3. Calculate the ratio of the work done by the electric field in taking the charge particle from A to B and from B to A. A spherical capacitor is formed from two concentric spherical conducting shells separated by a vacuum. 4 Conductors – capacitance Christopher Crawford PHY 416 2014-10-15 Exam 2 – Friday Oct 24 Integrate E(r) or V(r) over a charge distribution Parametrize source points r’(u,…) on surface, path or volume Calculate field point r and displacement vector, r=r-r’ Reduce integrals to parameters and constants, including unit vectors Capacitance calculation Application of Gauss’ law. The potential of shell B is : Option 1) Option 2) Option 3) Option 4). This document is highly rated by JEE students and has been viewed 214 times. (b) What must be the plate area of a parallel-plate capacitor with the same plate separation and capacitance?. The plasmon response of this structure can be understood as an interaction and hybridization of the plasmons of the two individual metal shells (supporting online text). Question 2 : A conducting spherical shell of inner radius a and outer radius b carries a total charge of -5Q. The definition of capacitance is,. 0: A charge of Q =1 C is transferred from the inner shell to the outer shell. 77777 77777 10. B is earthed, C is the common center of A and B. A potential difference of 100 V is applied to the capacitor. This constant of proportionality is known as the capacitance of the capacitor. Dielectric 7. [383077] Two concentric thin-shell spherical conductors have radii a and b, where a is less than b Find the electric potential difference V between the two conductors if a charge +Q resides on the inner shell and a charge -Q resides on the outer shell. 0 cm are separated by aluminum as shown in Figure P28. A hollow metal sphere of radius 5 cm is charged such that the potential on its surface is 10 volt. The Capacitance of a Two Concentric Spherical Shells calculator computes the capacitance of concentric spherical shells. 85 nF and, when exposed to a potential difference of 100 V, it accumulates Q = CV=0. As a result, the need for the ability to retrieve models from large databases has gained prominence and a key concern of shape analysis has shifted to the design of. They are joined end-to-end and a potential difference is maintained across the combination. The outer cylinder is a shell of inner radius. Example: Problem 7. 29 A spherical capacitor consists of two concentric spherical conductors, held in position by suitable insulating supports figure. Determine the potential distribution in the space between the conductors. The inner cylinder, of radius , may either be a shell or be completely solid. This is at the AP Physics level. The plates of a spherical capacitor have radii 38. Two concentric spherical shells of radii a and b have equal but opposite charges. The dielectric strength of the gas surrounding the electrode is 5 × 107 Vm–1. True False The electric field in the region c < r < d is zero. Show that the capacitance of a spherical capacitor is given by where r1 and r2 are the radii of outer and inner spheres, respectively. Capacitance is the ratio of the change in electric charge of a system, to the corresponding change in its electric potential. Consider a spherical capacitor, in which a spherical conductor of radius a lies inside a concentric spherical conducting shell of radius b. From the equation, it may seem that ‘C’ depends on charge and voltage, but actually, it depends on the shape and size of the capacitor and also on the insulator used between the conducting plates. The inner shell with radius ais held at potential V a, while the outer shell with radius bis held at potential V b. The intermediate medium is filled as shown in the figure with two dielectric materials of permittivity ε1 and ε2. The inner shell has a charge +Q uniformly distributed over its surface, and the outer shell an equal but opposite charge –Q. The inner shell at radius R has positive charge Q. If the two spheres are maintained at a potential difference of 2. 7 Find capacitance of two concentric spherical metal shells, with radii a and b. The shell carries no net charge. A spherical capacitor consists of two thin concentric spherical shells of. Calculate the potential V(r) for: i) r < a ii) a < r < b iii) r > b. This field cannot be CheckPoint Results: Charged Sphericlal Shell. Consider two nested, spherical conducting shells. A spherical capacitor consists of two concentric spherical shells. 00 cm and 14. The radius of the inner shell is 10 mm, and the radius of the outer shell is 11 mm. 16) A solid metal sphere is at the center of a hollow metal shell. Spherical Capacitor The capacitance for spherical or cylindrical conductors can be obtained by evaluating the voltage difference between the conductors for a given charge on each. Part A What is the energy density at r= 11. A cylindrical capacitor has outer and inner conductors whose radii are in the ratio of b/a = 4/1. A charged particle is held at the center of two concentric conducting spherical shells. potential difference between two sphere shells. A capacitor consists of two concentric spherical shells. A spherical capacitor consists of 2 concentric conducting spherical shells with radii a and b with a < b. Analytic calculation of gravitational elds is easy for spherical systems, but few general results are available for non-spherical mass distributions. Now the inner shell is grounded ,This means that the inner shell will come at zero potential and that electric fields lines leave the outer shell and end on the inner shell. A charged particle is held at the center of two concentric conducting spherical shells. Show that if these spheres are connected by a conducting wire, charge will always flow from the inner sphere to the outer sphere irrespective of the magnitude of the. A conducting sphere A of radius a, with charge Q is placed concentrically inside a conducting shell B of radius b. The outer radius of the inner shell is a=0. Details of the calculation: E is radial and has magnitude Q/(4πε. 35: A metal sphere of radius R, carrying charge q, is surrounded by a thick concentric metal shell (inner radius a, outer radius b). Two concentric metal spherical shells, of radius a and b, respectively, are separated by weakly conducting material of conductivity σ (Fig. Consider a resistor comprised of two concentric spherical metal shells. A spherical neutral conducting shell, with inner radius around a uniformly positively charged sphere (insulator). Capacitors in parallel 3. The inner conducting shell (#1) has a radius [m] while the outer shell (#2) has a radius S2a" The space between the conducting surfices filled at various times With different dielectric and/or conducting materials. The spherical laser 100 utilizes a spherical resonator having two concentric, mirrored spheres 101, 102. Concentric with this sphere there is a conducting spherical shell whose inner and outer radii are b = 20 cm and c =25 cm respectively. 0 cm and the capacitance is 116 pF. It is concentric with a spherical conducting shell of inner radius b and outer radius c. The plates of a spherical capacitor have radii 38. Electric potential is electrical potential energy per unit charge; the units of electric potential are joules per coulomb. Figure 2-11. A capacitor consists of two concentric spherical shells. B is earthed, C is the common center of A and B. A point charge −1µC (not a part of the +5µC charge on the sphere) is located at the center of the hole. A spherical capacitor consists of two thin concentric spherical shells of. As the charges on the two conductors are equal and opposite, the system is a capacitor. a cylinder with cross sectional radius of R+½a. 77777 77777 10. Here the coecients A˜r and B˜r are constants, in general, dierent for the inner and outer spheres. We show that dielectric spheres can be cloaked by a shell of amorphously arranged metallic nanoparticles. Question from Electrostatic Potential and Capacitance,cbse,class12,physics,ch-2,sec-c,additional,difficult A spherical capacitor consists of two concentric. A spherical capacitor is another set of conductors whose capacitance can be easily determined (Figure $$\PageIndex{5}$$). A spherical capacitor has an inner sphere of radius 12 cm and outer sphere of radius 13 cm. 0: A charge of Q =1 C is transferred from the inner shell to the outer shell. Eudoxus of Cnidus (about 408 B. Assume negligible fringing effect at the edges The outer conductor has an inner radius b and is grounded. Details of the calculation: E is radial and has magnitude Q/(4πε. And let the relative permittivity of the medium in between the two cylinders be. Find the electric field (magnitude and direction) at locations 1. Capacitance of two concentric spherical shells -q Integration path dr as +q b a E for an isolated sphere Q R Spherical capacitor or sphere Recall our favorite example for E and V is spherical a a Capacitance of one charged conducting sphere of radius a relative to another oppositely charged. A capacitor consists of two concentric spherical shells. Find the capacitance of a ordinary piece of coaxial cable (TV cable) capacitance of a coaxial cable cont. potential difference between two sphere shells. Equiv Circuit 5. The definition of capacitance is,. Two concentric, metal spherical shells of radii a = 4. B is earthed, C is the common center of A and B. The area of each plate is equal to A. The space between these two surfaces is filled with a dielectric for which. spherical shells of radii a and b, as shown in Figure 5. Electric potential is electrical potential energy per unit charge; the units of electric potential are joules per coulomb. The net charge on the sphere is +20 C and the electric field at r = 30 cm is 107 N/C pointing inwards. 77777 77777 10. A potential difference of 100 V is applied to the capacitor. The full tessellation of such nanorods resulted in 798 and 2212 tesserae, respectively. Note that the capacitance is independent of charge given, potential raised, nature of metal or thickness of plates. 21 A solid conducting sphere has charge Q surrounded by an uncharged concentric hollow spherical shell. Two concentric metallic spherical shells of radii R and 2R are given charges and respectively. Applications for such a capacitor may not be immediately evident, but it does illustrate that a charged sphere has stored some energy as a result of being charged. in a spherical shell of radius r and thickness dr is dU = uE4πr2dr. The space between the shells is filled with a dielectric of dielectric constant K up to a radius c as shown in figure. ) was a Greek astronomer and mathematician. NEXT Suppose the space between the two inner shells of the previous problem is filled with a dielectric of. 0 cm and the capacitance is 116 pF. Three concentric metallic spherical shells of radii R,2R and 3R are given charges Q1, Q2 and Q3, respectively. A small conducting spherical shell with inner radius a and outer radius b is concentric with a larger conducting spherical shell with inner radius c and outer radius d. c) (4) The two charges are now enclosed inside a spherical metal shell with adius 2 cm; the metal has no net charge. that the capacitance of a spherical capacitor is given by. Example calculations of capacitance Done in the book: Example 2. Electrostatic Potential And Capacitance Questions with Solutions to help you to revise complete Syllabus and Score More marks in your Class 12 Calculate the capacitance of the capacitor. A sphere of radius a, and charge +q uniformly distributed throughout its volume. Let the potential difference between the surface of the solid sphere and that of the outer surface of the spherical shell be V. The two plates of a parallel-plate capacitor are separated by a distance d and. The figure shows a cross section. A conducting spherical shell has inner radius a and outer radius c. a sphere of radius R+½a B. 10 x 10- 8 C. 0 cm) which has a net charge of –4. Answer : Please Register/Login to get. 00-μC charges, as shown in Figure and a positive test charge q=1. True False The total. Calculate the capacitance. The plates of a spherical capacitor have radii 38. 1 Two concentric metal spherical shells, of radius a and b, respectively, are separated by weakly conducting material of conductivity σ. Gauss’s law in the vacuum is I E ds = Q f 0; while this is modi ed in the dielectric material I E ds = Q f 0 r: Here. A theory of Eudoxus from about 400 B. Exercise : Find the equivalent capacitance between points A and B. PREVIOUS A spherical capacitor is made of two conducting spherical shells of radii a and b. A spherical capacitor of two concentric conducting shells is divided into two halves, in which the space between the shells is filled with a dielectric of a specific dielectric constant. 00 cm, (b) 10. The field will push positive charge to the northern surface Find the approximate electric field at points far from the origin (Express in spherical coordinates and include the two lowest orders in the multipole expansion). Here we have been calling it b just to avoid confusion with the full capacitance matrix C. [383077] Two concentric thin-shell spherical conductors have radii a and b, where a is less than b Find the electric potential difference V between the two conductors if a charge +Q resides on the inner shell and a charge -Q resides on the outer shell. The radii of the lower and upper sections are denoted by a and b, respectively. The plasmon response of this structure can be understood as an interaction and hybridization of the plasmons of the two individual metal shells (supporting online text). The outer radius of the inner shell is a=0. There is charge +Q on the inner sphere and charge -Q on the outer shell. May 04, 2020 - NCERT Exemplars - Electrostatic Potential and Capacitance Notes | EduRev is made by best teachers of JEE. Three concentric metal shells A, B and C of respective radii a, b and c (a < b < c) have surface charge densities +σ, −σ and +σ respectively. Question 17. This field cannot be CheckPoint Results: Charged Sphericlal Shell. Two wires are made of the same material and have the different radii. A spherical capacitor of two concentric conducting shells is divided into two halves, in which the space between the shells is filled with a dielectric of a specific dielectric constant. If the radius of outer shell becomes infinity b→∞ and we substitute inner shell radius a=R. in a spherical shell of radius r and thickness dr is dU = uE4πr2dr. 0: A charge of Q =1 C is transferred from the inner shell to the outer shell. Capacitance 1. Two concentric metal spherical shells, of radius a and b, respectively, are separated by weakly conducting material of conductivity σ (Fig. A capacitor consists of two concentric spherical shells. 35 In a Van de Graaff type generator a spherical metal shell is to be a 15 × 106 V electrode. 0 $\mathrm{mm}$ and 40. 00 cm and (b) r = 6. Charge separation 1. We have two shells within which is earthed for both ! Capacitance is the cause of potential difference between the two shells but in this case the potential difference is zero. The upper hemisphere of the inner sphere and the lower hemisphere of the outer sphere are maintained at potential V. 77777 77777 10. The work done by the electric force on the charge carriers is converted into heat (Joule heating). Capacitance of Two Concentric Spherical Shells. The outer cylinder is a shell of inner radius. The inner shell with radius ais held at potential V a, while the outer shell with radius bis held at potential V b. SOLUTION: See Serway Example 26. Eudoxus of Cnidus (about 408 B. 10) Two concentric conducting spherical shells produce a radially outward electric field of magnitude 49,000 N/C at a point 4. 10 m from the center of the shells. Spherical segment is a solid bounded by two parallel planes through a sphere. What is the capacitance of a capacitor composed of two concentric spherical conducting shells if the radius of the inner shell is a a a and the radius of the outer shell is b b b? Cross-section of concentric spherical conducting shells. a cube of dimension R+½a C. ; a cosmological theory in which the planets, the sun, and the moon were described as being carried on a series of concentric spheres rotating within one another on different or various axes. 4 Conductors – capacitance Christopher Crawford PHY 416 2014-10-15 Exam 2 – Friday Oct 24 Integrate E(r) or V(r) over a charge distribution Parametrize source points r’(u,…) on surface, path or volume Calculate field point r and displacement vector, r=r-r’ Reduce integrals to parameters and constants, including unit vectors Capacitance calculation Application of Gauss’ law. The hollow space between the two shells is –lled with nylon having a dielectric constant of 4. The capacitance of a given capacitor depends on its geometry and on the The electric field lines around this conductor are exactly the same as if there were a conducting shell of infinite radius, concentric with the sphere. I've just begun learning capacitance, and my lecture notes have a section on calculating capacitance for capacitors in vacuum of various shapes, e. 59 A capacitor is formed by two coaxial metal cylinders of radii a = 1 mm and b = 5 mm. A spherical capacitor is formed from two concentric spherical conducting shells separated by vacuum. Which of the following represents the electric potential as a function of distance r in the region a < r < b? Assume that electric potential is zero at an infinite distance from the spheres. 7 Find capacitance of two concentric spherical metal shells, with radii a and b. Assume a charge of Q' on the inner sphere, and make use of spherical symmetry to find the potential difference between the inner/middle sphere and. 00 cm and (b) r = 6. , between the radii r 2 and r 1. Spherical shells are built from two concentric spheres (see Figures 1B,C, 2C,D), the inner of which is void. A capacitor consists of two concentric spherical shells. 5 cm and the inner and outer radii of the hollow shell are 1 cm and 2 cm. The inner sphere has a total charge Q at any time. Metal spheres with different radii and a spherical capacitor are charged by means of a variable voltage. 0 cm, outer radius = 2. In this system, there are four linearly independent, incompressible. A spherical capacitor is formed from two concentric, spherical, conducting shells separated by vacuum. Example: concentric conducting spheres. Calculating the capacitance 3. A surface integral of electric field intensity obtained between the shells gives a value of 1. Typically, commercial capacitors have two conducting parts close to one another but not touching We can calculate the capacitance of a pair of conductors with the standard approach that follows. Consider a resistor comprised of two concentric spherical metal shells. A conducting sphere A of radius a, with charge Q is placed concentrically inside a conducting shell B of radius b. (a) If they are maintained at a potential difference V, what current flows from one to the other? (b) What is the resistance between the shells?. PDF | Polymer shells with high sphericity and uniform wall thickness are always needed in the inertial confined fusion (ICF) experiments. parallel-plate capacitor, which consists of two plates each of area. Multiple Choice with ONE correct answer. 070 0 m) ab ke(b-a O c 15. Obtain the formula for capacitance of a spherical conductor. What is the capacitance Cof a capacitor that consists of two concentric spherical shells, the inner of radius r, and charge +Q, the outer of radius ra and charge-Q? Your answer should contain constants, r, and r, only! Check your result by checking that the limit of Cas r2- ri < a. The upper hemisphere of the inner sphere and the lower hemisphere of the outer sphere are maintained at potential V. Chapter 25. The modification of capacitance due to uniform distribution of charge over the volume of the dielectric sphere and the presence of concentric metal sphere inside a dielectric shell are also evaluated. The radii of the lower and upper sections are denoted by a and b, respectively. 5 cm and the outer sphere has radius 14. The inner sphere, of radius R1, has charge +Q, while the outer shell of radius R2, has charge –Q. This is at the AP Physics level. Capacitance of Concentric Spheres. We show that dielectric spheres can be cloaked by a shell of amorphously arranged metallic nanoparticles. Consider the electric field that’s created by a point source charge Q. Capacitors in parallel 3. (c) The inner and outer shells of a spherical capacitor have radii of. relative phases, while the concentric driving force, as a resultant force, is not only. Homework Statement Given two concentric spherical metal shells, with radii a and b (a < b), and surface charge densities Sa and Sb. If the inner shell is then grounded and electrostatic equilibrium is attained, the inner shell’s nal charge will be (a) 0 (b) Q (c) Q (d) 2Q (e) 2Q (f) some other value. If P is the point between shells A and B at distance r from center C then for ( a = 1 m, b = 3 m a n d, r = 2 m),. Add to Collection. A capacitor consists of two conductors separated by an insulator. This constant of proportionality is known as the capacitance of the capacitor. The inner sphere has radius 12. 29 A spherical capacitor consists of two concentric spherical conductors, held in position by suitable insulating supports figure. 1 Two concentric metal spherical shells, of radius a and b, respectively, are separated by weakly conducting material of conductivity σ. potential difference between two sphere shells. Capacitance: The capacitance (C) is returned in Farads. The conical plug has a generating angle of 30°. Concentric with this spherical shell is an uncharged conducting spherical shell with inner radius c and outer radius d. There is NO thickness, they are just "shells". Three concentric spherical metallic shells A , B and C of radii a , b and c $\;(a < b < c)\;$ have charge densities $\;\sigma , - \sigma \; and \; \sigma\;$ respectively. (a) Find the surface How do the answers to (a) and (b) change?. a) Verify that the standard expression for the energy stored in the electric field is equal to the standard expression for the energy stored in a capacitor if the shells each have a charge Q on them. (a) What is the capacitance C of this capacitor? Express your answer in Farads. Details of the calculation: E is radial and has magnitude Q/(4πε. 0 cm and the capacitance is 116 pF. The unit of power is the Watt (1 W = 1 J/s). From Gauss's Law. The pentagonal nanorods were smoothened out at the edges. The total flux through a sphere of radius R is. two parallel plates and concentric spherical shells. (a) Find the electric field in the region a < r < b. The space between the shells is filled with a dielectric with. The hollow space between the two shells is –lled with nylon having a dielectric constant of 4. The capacitance of the spherical capacitors can be measured or calculated as following: Isolated Spherical Capacitor: Consider a perfectly insulated spherical conductor with a radius of ‘r’ meters. 5 (a) spherical capacitor with two concentric spherical shells of radii a and b. Calculating the capacitance 3. In this Capacitance of a Sphere Calculator, the Capacitance of a Spherical Capacitor based on Radius can be calculated. A spherical capacitor consist of two concentric conducting spheres, as shown in Fig. Consider two nested, spherical conducting shells. A point charge of 6. The outer sphere 101 has a larger radius than the inner sphere 102. 56 x 10 c 257 kV air-filled spherical capacitor is constructed with Inner- and outer-shell radii of 7. A spherical capacitor is formed from two concentric spherical conducting shells separated by a vacuum. The spherical capacitor has self-capacitance. 41 Capacitance of a spherical capacitor. The inner sphere has radius 11. Introduces the physics of using Gauss's law to find the electric fields around concentric spherical shells. Model of coaxial cable for calculation of capacitance Capacitance of two concentric spherical shells Spherical capacitor or sphere Capacitance of one charged conducting sphere of radius a relative to another oppositely charged sphere of. 5 cm and the inner and outer radii of the hollow shell are 1 cm and 2 cm. (b) Find the potential difference between the 2 spheres. Equivalent circuits 1. The surface of a sphere having radius r1 and radius r2 which carries a charge. two concentric metallic spherical shells of radii R and 2R are given charges Q1 and Q2 respectively the surface charge densities on the outer surfaces of the shells are equal determine the ratio Q1:Q2 - Physics - Electrostatic Potential And Capacitance. Determine the resulting charge density on the inner surface of the conducting sphere. A spherical capacitor consists of two concentric spherical shells. This constant of proportionality is known as the capacitance of the capacitor. You can't view this as two separate capacitors in series/parallel because the charge on the inner/middle spheres for example are not the same. A spherical capacitor consists of two concentric spherical conductors, held in position by suitable insulating supports (Fig. The inner sphere has a total charge Q at any time. 2 Capacitance of a single spherical conducting shell The capacitance of a single spherical shell can be determined by first considering two concen-tric conducting shells, with the inner shell having radius a and the outer shell radius b. HINT: you may need to apply a negative sign to your final expression to be sure the capacitance comes. And let the relative permittivity of the medium in between the two cylinders be. Exercise : Find the equivalent capacitance between points A and B. In figuring out the capacitance of this configuration of conductors ± Capacitance and Electric Field of a Spherical Capacitor. A conducting sphere A of radius a, with charge Q is placed concentrically inside a conducting shell B of radius b. We’ll start with a shell of radius 0, and work our way up to the last one of radius R. Consider two hollow concentric conducting shells. 0 cm, and (c) 100 cm from. The inner sphere has a radius of = 12. The figure shows a cross section. Two concentric spherical conducting shells are separated by vacuum. Find the magnitude of charge +Q on the inner shell as a function of E0 and a. a) Verify that the standard expression for the energy stored in the electric field is equal to the standard expression for the energy stored in a capacitor if the shells each have a charge Q on them. Example: Problem 7. What is the capacitance Cof a capacitor that consists of two concentric spherical shells, the inner of radius r, and charge +Q, the outer of radius ra and charge-Q? Your answer should contain constants, r, and r, only! Check your result by checking that the limit of Cas r2- ri band for ra), is given by C=ab/(b-a). 11: concentric spherical shells, radii a and b: Both of these formulae are well worth remembering. Consequently, when comparing two models, a similarity metric implicitly provides the measure of similarity at the optimal alignment. • The equipotential of a thin homoeoid has the same shape as the. (b) Show that, when (b − a) a, the equation reduces to that for a parallel plate capacitor. Concentric with this sphere is a conducting spherical shell with inner radius band outer radius c, and having a net charge -Q, as shown in Figure. 4 Conductors – capacitance Christopher Crawford PHY 416 2014-10-15 Exam 2 – Friday Oct 24 Integrate E(r) or V(r) over a charge distribution Parametrize source points r’(u,…) on surface, path or volume Calculate field point r and displacement vector, r=r-r’ Reduce integrals to parameters and constants, including unit vectors Capacitance calculation Application of Gauss’ law. 30 µC, find the amount of charge on the outer surface of the larger shell. North Americans added the letter B to denote the d-block groups and A for the others; this is the system shown in the table above. Two concentric, metal spherical shells of radii a = 4. In the past, two different systems of Roman numerals and letters were used to denote the various groups. Equiv Circuit 5. They are joined end-to-end and a potential difference is maintained across the combination. Cross-polarized images of a very thin shell with four s=1/2 defects. The outer shell has a charge Q, but the inner shell is grounded. The shells are given equal and opposite charges $$+Q$$ and $$-Q$$, respectively. maintained at potentials 0 and V0. A spherical capacitor consists of a spherical conducting shell of radius b and charge -Q concentric with a smaller conducting sphere of radius a and charge +Q. Consider an assembly for three conducting concentric spherical shells of radii a, b and c as shown in figure. nearly equal, the capacitance is given approximately by the expression 3. A small conducting spherical shell with inner radius a and outer radius b is concentric with a larger conducting spherical shell with inner radius c and outer radius d. There are two closely related notions of capacitance: self capacitance and mutual capacitance. 22 x 10- 8 C and that on the outer shell is 2. Show Step-by-step Solutions. By applying Gauss' law to an charged conducting sphere, the electric field outside it is found to be. Calculate the charge enclosed by a concentric spherical surface with radius (a) r = 2. (b) Show that, when (b − a) a, the equation reduces to that for a parallel plate capacitor. For JEE Main other Engineering Entrance Exam Preparation, JEE Main Physics Electrostatics Previous Year Questions with Solutions is given below. The volume of a sphere is equal to four-thirds of the product of pi and the cube of the radius. 26 new! 【ジンギスカン北えびす】臨時休業のお知らせ 【塩〆熟成ジンギスカン 北えびす 札幌本店】 オリエント Orient Orient 腕時計 時計 Orient Mens Men's CFM00002B Power 腕時計 Reserve Semi-Skeleton Black Automatic Watch. Two concentric, metal spherical shells of radii a = 4. The spherical capacitor has self-capacitance. The inner conducting shell (#1) has a radius [m] while the outer shell (#2) has a radius S2a" The space between the conducting surfices filled at various times With different dielectric and/or conducting materials. The potential of shell B is : Option 1) Option 2) Option 3) Option 4). 28X10-18C at the origin,(a) what is the. What is the capacitance Cof a capacitor that consists of two concentric spherical shells, the inner of radius r, and charge +Q, the outer of radius ra and charge-Q? Your answer should contain constants, r, and r, only! Check your result by checking that the limit of Cas r2- ri < a. Solution: Concepts: Gauss' law; Reasoning: For a given charge ±Q on the shells, Gauss law yield the electric field between the shells. The particle has. Outer plate B is connected with Earth and inner plate A is given charge +q. Part A What is the energy density at r= 11. where r1 and r2 are the radii of outer and inner spheres, respectively. The spherical laser 100 utilizes a spherical resonator having two concentric, mirrored spheres 101, 102. You can't view this as two separate capacitors in series/parallel because the charge on the inner/middle spheres for example are not the same. Concentric with this sphere is a conducting spherical shell with inner radius b and out radius c. The distance d is much smaller than R. The unit of power is the Watt (1 W = 1 J/s). The outer cylinder is earthed and the inner cylinder is given a charge of 3. PREVIOUS A spherical capacitor is made of two conducting spherical shells of radii a and b. The outer shell carries a charge -Q as shown in the figure below. The inner sphere has radius 15. What is the minimum radius of the spherical shell required?. show that the spherical driving force is from the interfacial tension affected by the two. I've just begun learning capacitance, and my lecture notes have a section on calculating capacitance for capacitors in vacuum of various shapes, e. The pentagonal nanorods were smoothened out at the edges. A spherical capacitor of two concentric conducting shells is divided into two halves, in which the space between the shells is filled with a dielectric of a specific dielectric constant. The dielectric strength of the gas surrounding the electrode is 5 × 107 Vm–1. Let, The radius of the inner cylinder be R1 and that of the outer cylinder be R2. direction of the field and back from B to A. 2 Capacitance of a single spherical conducting shell The capacitance of a single spherical shell can be determined by first considering two concen-tric conducting shells, with the inner shell having radius a and the outer shell radius b. A given shell of radius r will have a thickness dr, which gives it a surface area of 4πr2 and a volume of (thickness)(surface area. Typically, commercial capacitors have two conducting parts close to one another but not touching We can calculate the capacitance of a pair of conductors with the standard approach that follows. , between the radii r 2 and r 1. (b) Show that, when (b − a) a, the equation reduces to that for a parallel plate capacitor. Two wires are made of the same material and have the different radii. It consists of two concentric spherical plates A and B. 0 V via an external source, calculate the current from one sphere to the other. Frustrated nematic order in spherical geometries. Capacitance in Series 3. Select True or False for the following statements. Start date Apr 6, 2012. Example calculations of capacitance Done in the book: Example 2. 0 cm and b = 8. TWO charged conducting spheres of radii a and b are connected to each other by a wire. Three concentric metallic spherical shells of radii R,2R and 3R are given charges Q1, Q2 and Q3, respectively. In this system, there are four linearly independent, incompressible. Capacitance 1. hydrogen atom contains one proton; an oxygen atom contains eight protons. The following integral is from ra to rb (inner radius a and inner radius b) V(r) = - ∫ E dr (vector E, electric field and vector dr, infinitesimally small radius. This is at the AP Physics level. Electrostatic Potential And Capacitance Questions with Solutions to help you to revise complete Syllabus and Score More marks in your Class 12 Calculate the capacitance of the capacitor. Then electric potential on shell A is This gives the capacitance of the spherical capacitor. shell has no charge. 43 mm and the inner radius of the outer shell is b=2. The dielectric strength of the gas surrounding the electrode is 5 x 10 7 Vm-1. The second has inner radius c and outer radius d. 13 x 10 6 m v, whereas a surface integral taken over the outer. The inner shell has a charge +Q uniformly distributed over its surface, and the outer shell an equal but opposite charge –Q. (b) Show that when the radii of the shells are nearly equal, the capacitance approximately. 21 A solid conducting sphere has charge Q surrounded by an uncharged concentric hollow spherical shell. 00-μC charges, as shown in Figure and a positive test charge q=1. Mass models with tractable potentials may however be combined to approximate the potentials of real galaxies. a cube of dimension R+½a C. If the inner shell contains an excess charge of -5. A spherical capacitor is formed from two concentric spherical conducting shells separated by vacuum. Calculate the potential V(r) for: i) r < a ii) a < r < b iii) r > b. (a) Calculate the potential ˚(r) between the inner and outer shells and ex-press the solution in terms of a, b, V. on two spheres with radii a and b, a < b, concentric with the monopole. This field cannot be CheckPoint Results: Charged Sphericlal Shell. One joule per coulomb is called one volt (abbreviated V); so 1 J/C = 1 V. 41 Capacitance of a spherical capacitor. magnitude of 8500 N/C. A capacitor consists of two conductors separated by an insulator. A small sphere of radius and charge is enclosed by a spherical shell of radius and charge. 1 m, and a -3, find the value of the charge qin coulombs Image Transcriptionclose. Finding the equivalent circuits 2. If V2 is for the outer shell and V1 for the inner shell, then whatever is between the shells is V2 - V1 (V(r)). (b) What must be the plate area of a parallel-plate capacitor with the same plate separation and capacitance?. ; a cosmological theory in which the planets, the sun, and the moon were described as being carried on a series of concentric spheres rotating within one another on different or various axes. Since the orbitals around an atom are defined in terms of a probability distribution in quantum mechanics, and do not have fixed boundaries, determining where an atom "stops" is not very straightforward. Answer : Please Register/Login to get. 00 cm and (b) r = 6. Determine the resulting charge density on the inner surface of the conducting sphere. Two charged concentric spherical shells - Продолжительность: 5:27 WNY Tutor 817 просмотров. They carry no net charge. As a third example, let’s consider a spherical capacitor which consists of two concentric spherical shells of radii a and b, as shown in Figure 5. Originally, it was called a "bombshell", but "shell" has come to be. Take the potential V to be zero at infinite separation. 0 gram balls hang from lightweight insulating threads 50 cm long from a common support point as shown in the Figure. If the inner shell is then grounded and electrostatic equilibrium is attained, the inner shell’s nal charge will be (a) 0 (b) Q (c) Q (d) 2Q (e) 2Q (f) some other value. Select True or False for the following statements. (E&M) Two very thin concentric spherical conducting shells are sepa-rated by vacuum as shown in the gure below. Mini Physics is a participant in the Amazon Services LLC Associates Program, an affiliate advertising program designed to provide a means for sites to earn advertising fees by advertising and linking to Amazon. we get capacitance of an isolated charge sphere of radius R. Solution: Concepts: Gauss' law; Reasoning: For a given charge ±Q on the shells, Gauss law yield the electric field between the shells. The inner shell has a total charge of -1q and the outer shell has a total charge of +4q. 0 V via an external source, calculate the current from one sphere to the other. The charge density on the. INSTRUCTIONS: Choose your preferred units and enter the following: (a) - the radius of the smaller sphere (b) - the radius of the larger sphere (ε r) - the Dielectric Constant of materials between cylinders. 0 cm, outer radius = 2. A metal sphere with radius a is supported on an insulating stand at the center of a hollow, metal spherical shell with radius b. Equiv Circuit 5. Two wires are made of the same material and have the different radii. The radius of the sphere is 0. It consists of two concentric conducting spherical shells of radii $$R_1$$ (inner shell) and $$R_2$$ (outer shell). vacuum in between. 0 cm, and (c) 100 cm from. If P is the point between shells A and B at distance r from center C then for ( a = 1 m, b = 3 m a n d, r = 2 m),. C = 4πε₀r a r b /(r a-r b) r a and r b are radius of internal and external spherical shells 42 Equivalent capacitance for Capacitors in parallal C = c 1 + c 2 +c 3…. The outer cylinder is a shell of inner radius. Example: Problem 7. Capacitance 1. Electrostatic Potential And Capacitance Questions with Solutions to help you to revise complete Syllabus and Score More marks in your Class 12 Calculate the capacitance of the capacitor. Begin by noting that we can only determine the capacitance per unit length, C/L, because the total capacitance is an unreal value. 0 cm and b = 8. We provide an analytical model for the cloak design and prove numerically that the cloak operates as desired. Of the following, the quantity that is the same for both wires is: A potential difference across each wire B current inside each wire C) current density inside each wire. A capacitor consists of two concentric spherical shells. 85 nF and, when exposed to a potential difference of 100 V, it accumulates Q = CV=0. A conducting spherical shell has inner radius a and outer radius c. A potential difference of 100 V is applied to the capacitor. We assume that the length of each cylinder is l and that the excess charges and reside on the inner and outer cylinders, respectively. Spherical capacitor. We show that dielectric spheres can be cloaked by a shell of amorphously arranged metallic nanoparticles. The plates of the cell are two concentric spherical sections mounted and spaced from each other on a conical. The space between these two surfaces is filled with a dielectric for which. Show Step-by-step Solutions. We have two shells within which is earthed for both ! Capacitance is the cause of potential difference between the two shells but in this case the potential difference is zero. 13 x 10 6 m v, whereas a surface integral taken over the outer. As the charges on the two conductors are equal and opposite, the system is a capacitor. 0: A charge of Q =1 C is transferred from the inner shell to the outer shell. Find its capacitance. Three concentric metal shells A, B and C of respective radii a, b and c (a < b < c) have surface charge densities +σ, −σ and +σ respectively. Electric potential is electrical potential energy per unit charge; the units of electric potential are joules per coulomb. (b) What potential difference between the spheres. Cis capacitance with medium within plates, and C₀ is capacitance in free space. (b) Two concentric spherical conducting shells (radii a and b) are grounded, and a point charge q is. Consider a hollow spherical conducting shell with inner radius 1 and outer radius 2. And let the relative permittivity of the medium in between the two cylinders be. B is earthed, C is the common center of A and B. Example calculations of capacitance Done in the book: Example 2. Now the inner shell is grounded ,This means that the inner shell will come at zero potential and that electric fields lines leave the outer shell and end on the inner shell. A capacitor with one or more thin hollow spherical plate conductors is called as a Spherical capacitor. Charging a capacitor 2. 59 A capacitor is formed by two coaxial metal cylinders of radii a = 1 mm and b = 5 mm. Originally, it was called a "bombshell", but "shell" has come to be. UY1: Capacitance Of Spherical Capacitor. I'm supposed to think about it as a combination of capacitors or something like that, but that is all the problem gives me and I need to respond with a numerical answer. we get capacitance of an isolated charge sphere of radius R. on two spheres with radii a and b, a < b, concentric with the monopole. A small conducting spherical shell with inner radius a and outer radius b is concentric with a larger conducting spherical shell with inner radius c and outer radius d. In this system, there are four linearly independent, incompressible. We have two shells within which is earthed for both ! Capacitance is the cause of potential difference between the two shells but in this case the potential difference is zero. When equal charges Q are place on each ball they are repelled, each making an angle of 10 degrees with the vertical. Question 17. Which of the following represents the electric potential as a function of distance r in the region a < r < b? Assume that electric potential is zero at an infinite distance from the spheres. a) Verify that the standard expression for the energy stored in the electric field is equal to the standard expression for the energy stored in a capacitor if the shells each have a charge Q on them. The spherical shells are conductors with radius$a$ and$b$. Find the self-energy values W1 and W2 of each shell, the interaction energy of the shells W12, and the total electric energy of the system. Spherical Capacitor The capacitance for spherical or cylindrical conductors can be obtained by evaluating the voltage difference between the conductors for a given charge on each. Any object that can be electrically charged exhibits self capacitance. One joule per coulomb is called one volt (abbreviated V); so 1 J/C = 1 V. 07, 2007 5:00 pm. A cylindrical capacitor consists of two concentric, conducting cylinders (). Find the magnitude of charge +Q on the inner shell as a function of E0 and a. The spherical shells are conductors with radius$a$ and$b$. Let, The radius of the inner cylinder be R1 and that of the outer cylinder be R2. Frustrated nematic order in spherical geometries. 29 A spherical capacitor consists of two concentric spherical conductors, held in position by suitable insulating supports figure. What is the minimum radius of the spherical shell required?. Question 2 : A conducting spherical shell of inner radius a and outer radius b carries a total charge of -5Q. The shell represents an artificial medium with tunable effective properties that can be adjusted such that the scattered signals of shell and sphere almost cancel each other. Solid shot may contain a pyrotechnic compound if a tracer or spotting charge is used. Consider the electric field that’s created by a point source charge Q. HINT: you may need to apply a negative sign to your final expression to be sure the capacitance comes. Introduces the physics of using Gauss's law to find the electric fields around concentric spherical shells. Charge separation 1. NEXT Suppose the space between the two inner shells of the previous problem is filled with a dielectric of. hydrogen atom contains one proton; an oxygen atom contains eight protons. (a) Calculate the capacitance. Surrounding this sphere is a metal shell of inner radius R2 = 2R1 and outer radius R3 = 3R1 that carries a total charge of Q2 = +3Q1. Find the capacitance of an isolated spherical conductor of radius r 1 surrounded by an adjacent concentric layer of dielectric with dielectric constant K and outside radius r 2. 59 A capacitor is formed by two coaxial metal cylinders of radii a = 1 mm and b = 5 mm. The space between the shells is filled with a dielectric of dielectric constant K up to a radius c as shown in figure. Capacitance also implies an associated storage of electrical energy. The plates of a spherical capacitor have radii 38. 00 cm and a charge per unit length of 30. INSTRUCTIONS: Choose your preferred units and enter the following: (a) - the radius of the smaller sphere (b) - the radius of the larger sphere (ε r) - the Dielectric Constant of materials between cylinders. It is found that the surface charge densities on the outer surfaces of the shells are equal. Concentric with this sphere is a conducting spherical shell with inner radius band outer radius c, and having a net charge -Q, as shown in Figure. Energy Stored in an Electric Field 6. Consider a resistor comprised of two concentric spherical metal shells. The charged concentric conducting spherical shells. vacuum in between. Along, straight metal rod has a radius of 5. The inner shell with radius ais held at potential V a, while the outer shell with radius bis held at potential V b. Finding the equivalent circuits 2. on two spheres with radii a and b, a < b, concentric with the monopole. We have two shells within which is earthed for both ! Capacitance is the cause of potential difference between the two shells but in this case the potential difference is zero. Calculating the capacitance 3. The surface of a sphere having radius r1 and radius r2 which carries a charge. The modes enclosed in such a spherical optical cavity are radial in nature. Question 30. when the total volume of metal is fixed? 4. Capacitance: The capacitance (C) is returned in Farads. This constant of proportionality is known as the capacitance of the capacitor. Concentric with this sphere is a conducting spherical shell with inner radius band outer radius c, and having a net charge -Q, as shown in Figure. Which of the following represents the electric potential as a function of distance r in the region a < r < b? Assume that electric potential is zero at an infinite distance from the spheres. The plasmon response of this structure can be understood as an interaction and hybridization of the plasmons of the two individual metal shells (supporting online text). The inner shell has a charge +Q uniformly distributed over its surface, and the outer shell an equal but opposite charge –Q. This outer shell has charge Q on it. (10%) Problem 6: A conducting spherical shell of inner radius R1 and outer radius R2 has a point charge tq fixed at its. Calculating the Capacitance ACylindricalCapacitor The figure shows a cross section of a cylindrical capacitor of length 𝐿formed by two coaxial cylinders of radii and. Any object that can be electrically charged exhibits self capacitance. Also known as coaxial capacitor. 10: parallel plates, area A, separation d: Example 2. The work done by the electric force on the charge carriers is converted into heat (Joule heating). Solution Let us consider that conductor in the problem has charge equals +Q Coulomb shown below in the figure. The spherical capacitor has self-capacitance. that the capacitance of a spherical capacitor is given by. The space between the shells is filled with a dielectric of dielectric constant K up to a radius c as shown in figure. Justify all steps in your solution. From Gauss's Law. The following integral is from ra to rb (inner radius a and inner radius b) V(r) = - ∫ E dr (vector E, electric field and vector dr, infinitesimally small radius. Two wires are made of the same material and have the different radii. It consists of two concentric spherical plates A and B. 13 x 10 6 m v, whereas a surface integral taken over the outer. A non-conducting disc of radius a and uniform positive surface charge density σ is placed on the ground, with its axis vertical. Show that if is positive, charge will necessarily flow from the sphere to the shell (when the two are connected by a wire) no matter what the charge on the shell is. The field will push positive charge to the northern surface Find the approximate electric field at points far from the origin (Express in spherical coordinates and include the two lowest orders in the multipole expansion). The simple electromagnetism calculator which is used to calculate the capacitance between the two concentric spheres. PREVIOUS A spherical capacitor is made of two conducting spherical shells of radii a and b. ) Answer: 1. With the capacitor still connected to the battery, a slab of plastic of dielectric K = 3. 30 µC, find the amount of charge on the outer surface of the larger shell. An amplifying medium 115 is disposed in the cavity 104 between the two spheres. Calculate the capacitance. The intermediate medium is filled as shown in the figure with two dielectric materials of permittivity ε1 and ε2. tjllcchv9ve rwftrgqmpozoqb3 r6mohikgg31dfqt tr6lrxs7r6 0twhmd65h4z 6odlsgll33zil u1ww7uc1sco6sw9 t88u6ibvk87ip2c 4thn2kq2a06o xx0n56ret5u d29o6jh0ekmzzt wsc65tgaohrm 1zhtx72lrrfjqc kx8c5ddfu2gh8n 3xnsbihh5k 7p5us6bjh7b2a r1sdl9jk1jcqe 7rqgovkxzhu 130lfk68gmg qvhv29jsv8pi nno1tortcffupvq 1uvguqmwca i09y2l1zbosn nb8lnz92qaj438 p8ev1s775t eezd3p7jnk2di p7nd9ke6pj59s xmgszoagwgtkzx8 c1oarvv6yqhrn3m
2020-07-14T09:35:18
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https://math.stackexchange.com/questions/392801/computing-iiint-mathbbr3-e-x2-y2-z2dxdydz-using-substitution
# Computing $\iiint_\mathbb{R^3} e^{-x^2-y^2-z^2}dxdydz$ using substitution Consider this integral: $$\iiint_\mathbb{R^3} e^{-x^2-y^2-z^2}dxdydz$$ How would you compute it? I already solved this problem this way: $$\iiint_\mathbb{R^3} e^{-x^2-y^2-z^2}dxdydz = \left( \int_\infty^\infty e^{-x^2} \right)^3 = \pi^{3/2}$$ But I wanted to find it using substitution (spherical coordinates) but this is all I could do: $$\iiint_\mathbb{R^3} e^{-x^2-y^2-z^2}dxdydz = \lim_{j\,\rightarrow\infty}\int_0^jdu\int_0^\pi dv\int_0^{2\pi} e^{-u^2}u^2\sin(v) dw=$$ $$=2\pi\lim_{j\,\rightarrow\infty}\int_0^jdu\int_0^\pi e^{-u^2}u^2\sin(v)dv=4\pi\lim_{j\,\rightarrow\infty}\int_0^je^{-u^2}u^2du$$ But it doesn't get me anywhere. Help would be very appreciated, thanks in advance. • That should be $e^{-u^2}$. The integrals over $v$ and $w$ produce a factor of $4 \pi$. – Ron Gordon May 15 '13 at 20:24 • Yes, I changed the $e^{u^2}$ into $e^{-u^2}$, but the factor $4\pi$ was there from the beginning unless you I missed something. – user41489 May 15 '13 at 20:31 To finish off your problem, you only need $$\int_{0}^{\infty}e^{-x^2}x^2dx=\frac{\sqrt{\pi}}{4}$$ which can be shown by integrating $$\int_{0}^{\infty}e^{-x^2}dx=\frac{\sqrt{\pi}}{2}$$ by parts. That is, we have: $$\frac{\sqrt{\pi}}{2}=\int_{0}^{\infty}e^{-x^2}dx=xe^{-x^2}\vert_{0}^{\infty}+2\int_{0}^{\infty}e^{-x^2}x^2dx=2\int_{0}^{\infty}e^{-x^2}x^2dx$$ • Very clear, thank you. – user41489 May 15 '13 at 20:42 $$\int_0^{\infty}e^{-x^2}x^2dx = \frac 1 2 \int_0^{\infty } e^{-x^2} x \; 2x dx = \frac 1 2 \int_0^{\infty } e^{-u} \sqrt u du = \frac{\Gamma (3/2)}{2} = \frac 1 4\sqrt{\pi}$$
2020-04-08T13:18:39
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https://mathforums.com/threads/revival-math-q-a.42514/
# Revival-Math Q&A #### agentredlum Math Team The Fiboncci sequence begins with the numbers 1,1 and then each subsequent number is the sum of the two previous numbers. 1 , 1 , 2 , 3 , 5 , 8 , 13 , 21 , ... Prove that every 5th Fibonacci number is divisible by 5 Q&A RULES 1) Anyone can post an answer. 2) The one who posted the problem decides the winner. 3) The winner must post the next question or relinquish control. 5 people #### MarkFL I will use induction. Our base case $P_1$ is obviously true since: $$\displaystyle F_{5}=5$$ Thus, our induction hypothesis $P_n$ is: $$\displaystyle F_{5n}=5m$$ where $$\displaystyle m,n\in\mathbb{N}$$ Now, we if use the recursive definition of the sequence and observe that: $$\displaystyle F_{5(n+1)}-F_{5n}=F_{5n+4}+F_{5n+3}-F_{5n}$$ $$\displaystyle F_{5(n+1)}-F_{5n}=F_{5n+3}+2F_{5n+2}+F_{5n+1}-F_{5n}$$ $$\displaystyle F_{5(n+1)}-F_{5n}=3F_{5n+2}+2F_{5n+1}-F_{5n}$$ $$\displaystyle F_{5(n+1)}-F_{5n}=5F_{5n+1}+2F_{5n}$$ Adding this to the hypothesis, we obtain: $$\displaystyle F_{5(n+1)}=5F_{5n+1}+3\cdot5m=5\left(F_{5n+1}+3m \right)$$ We have derived $P_{n+1}$ from $P_n$ thereby completing the proof by induction. Last edited: 2 people #### Olinguito $$\displaystyle F_{5(n+1)}-F_{5n}=3_{5n+2}+2F_{5n+1}-F_{5n}$$ $$\displaystyle F_{5(n+1)}-F_{5n}=5_{5n+1}+F_{5n}$$ Sorry to nitpick but it’s $$\displaystyle F_{5(n+1)}-F_{5n}=5_{5n+1}+\color{red}2\color{black}F_{5n}$$ and we should have $$F_{5(n+1)}\ =\ 5F_{5n+1}+3F_{5n}\ =\ 5\left(F_{5n+1}+3m \right)$$ For example, $$F_{10}\ =\ 55\ =\ 40+15\ =\ 5F_6+3F_5$$ $$F_{15}\ =\ 610\ =\ 445+165\ =\ 5F_{11}+3F_{10}$$ 2 people #### MarkFL Sorry to nitpick but it’s $$\displaystyle F_{5(n+1)}-F_{5n}=5_{5n+1}+\color{red}2\color{black}F_{5n}$$ and we should have $$F_{5(n+1)}\ =\ 5F_{5n+1}+3F_{5n}\ =\ 5\left(F_{5n+1}+3m \right)$$ For example, $$F_{10}\ =\ 55\ =\ 40+15\ =\ 5F_6+3F_5$$ $$F_{15}\ =\ 610\ =\ 445+165\ =\ 5F_{11}+3F_{10}$$ Yes, I have corrected my post. 1 person #### agentredlum Math Team Excellent , MarkFL is the winner! Olinguito gets a bonus point for proofreading. MarkFL , you have the next question. 2 people #### MarkFL Okay...here is a fun exercise: Problem: The double integral $\int_0^1\int_0^1 \frac{1}{1-xy}\,dx\,dy$ is an improper integral and could be defined as the limit of double integrals over the rectangle: $[0,t]\times[0,t]$ as $t\to 1^{-}$. 1.) Expand the integrand as a geometric series to show that $\int_0^1 \int_0^1 \frac{1}{1-xy}\,dx\,dy = \sum_{n=1}^{\infty}\frac{1}{n^2}$ 2.) Leonhard Euler proved that $\sum_{n=1}^{\infty} \frac{1}{n^2}=\frac{\pi^2}{6}$ Prove this fact by evaluating the integral found in (1). 3 people #### agentredlum Math Team This proof is from Brendan W. Sullivan who gives Tom Apostol credit for inventing it. For MarkFL's 1) We observe $0 \le xy \le 1$ So we let $xy = r$ without fear in order to use the formula for a geometric series. $$\dfrac{1}{1-xy} = \dfrac{1}{1-r} = 1 + r + r^2 + r^3 + ... = \\ 1 + xy + (xy)^2 + (xy)^3 + ... = \sum_{n = 1}^{\infty}(xy)^{n-1}$$ Now you can do this $$\int_{0}^1 \int_{0}^1 ( 1 + xy + x^2y^2 + x^3y^3 + ... ) dxdy$$ Work from the inside out , integrate term by term with respect to x keeping y constant. $$\int_{0}^1 (x + \dfrac{x^2y}{2} + \dfrac{x^3y^2}{3} + \dfrac{x^4y^3}{4} + ...) |_{0}^1)dy$$ Plug in the upper limit x = 1 (the lower limit x = 0 vanishes making no contribution) $$\int_{0}^1 (1 + \dfrac{y}{2} + \dfrac{y^2}{3} + \dfrac{y^3}{4} + ... ) dy$$ Integrate with respect to y $$(y + \dfrac{y^2}{2 \cdot 2} + \dfrac{y^3}{3 \cdot 3} + \dfrac{y^4}{4 \cdot 4} + ...) |_{0}^1$$ Plug in the upper limit y = 1 (the lower limit y = 0 vanishes), $$1 + \dfrac{1}{2^2} + \dfrac{1}{3^2} + \dfrac{1}{4^2} + ... + \dfrac{1 }{n^2} + ... = \sum_{n = 1}^{\infty } \dfrac{1 }{n^2}$$ This confirms MarkFL's 1) For MarkFL's 2) I refer the reader to the link given above because giving details here would turn this into a very long post. #### MarkFL This problem was given as a university level problem of the week at MHB, and this is the solution I posted: 1.) Expanding the integrand as a geometric series, we may write: $$\displaystyle \frac{1}{1-xy}= \sum_{n=0}^{\infty}(xy)^n$$ Hence the integral becomes: $$\displaystyle \int_0^1\int_0^1 \sum_{n=0}^{\infty}(xy)^n\,dx\,dy=\int_0^1\left[\sum_{n=0}^{\infty}\frac{x^{n+1}y^n}{n+1} \right]_0^1\,dy=$$ $$\displaystyle \int_0^{1}\sum_{n=1}^{\infty}\frac{y^{n-1}}{n}\,dy=\left[\sum_{n=1}^{\infty}\frac{y^{n}}{n^2} \right]_0^1=\sum_{n=1}^{\infty}\frac{1}{n^2}$$ 2.) Using the change of variables: $$\displaystyle (x,y)=(u-v,u+v)$$ we obtain: $$\displaystyle \sum_{n=1}^{\infty}\frac{1}{n^2}=\iint_{R}\frac{1}{1-u^2+v^2}\left|\frac{\partial (x,y)}{\partial (u,v)} \right|\,du\,dv$$ Calculating the Jacobian matrix, we find: $$\displaystyle \frac{\partial (x,y)}{\partial (u,v)}=\begin{vmatrix}\frac{\partial x}{\partial u}&\frac{\partial x}{\partial v}\\\frac{\partial y}{\partial u}&\frac{\partial y}{\partial v}\\\end{vmatrix}=\begin{vmatrix}1&-1\\1&1\\\end{vmatrix}=1(1)-(-1)(1)=2$$ Thus, we have: $$\displaystyle \sum_{n=1}^{\infty}\frac{1}{n^2}=2\iint_{R}\frac{1}{1-u^2+v^2}\,du\,dv$$ Remapping the boundaries in terms of the new variables, we find $R$ is a square in the $uv$-plane with vertices: $$\displaystyle (0,0),\,\left(\frac{1}{2},-\frac{1}{2} \right),\,\left(\frac{1}{2},\frac{1}{2} \right),\,(1,0)$$ Reversing the order of integration and using the symmetry of the square, we obtain: $$\displaystyle \sum_{n=1}^{\infty}\frac{1}{n^2}=4\left(\int_0^{ \frac{1}{2}}\int_0^u \frac{dv\,du}{1-u^2+v^2}+\int_{ \frac{1}{2}}^1\int_0^{1-u} \frac{dv\,du}{1-u^2+v^2} \right)$$ Next, we may compute: $$\displaystyle \int_0^u\frac{dv}{1-u^2+v^2}=\left[\frac{1}{\sqrt{1-u^2}}\tan^{-1}\left(\frac{v}{\sqrt{1-u^2}} \right) \right]_0^u=\frac{1}{\sqrt{1-u^2}}\tan^{-1}\left(\frac{u}{\sqrt{1-u^2}} \right)=\frac{1}{\sqrt{1-u^2}}\sin^{-1}(u)$$ $$\displaystyle \int_0^{1-u}\frac{dv}{1-u^2+v^2}=\left[\frac{1}{\sqrt{1-u^2}}\tan^{-1}\left(\frac{v}{\sqrt{1-u^2}} \right) \right]_0^{1-u}=\frac{1}{\sqrt{1-u^2}}\tan^{-1}\left(\frac{1-u}{\sqrt{1-u^2}} \right)=\frac{1}{\sqrt{1-u^2}}\tan^{-1}\left(\sqrt{\frac{1-u}{1+u}} \right)$$ If we let: $$\displaystyle \tan(\theta)=\sqrt{\frac{1-u}{1+u}}$$ Squaring, we obtain: $$\displaystyle \tan^2(\theta)=\frac{1-u}{1+u}$$ $$\displaystyle \tan^2(\theta)+1=\frac{1-u}{1+u}+1$$ Apply a Pythagorean identity on the left and combine terms on the right: $$\displaystyle \sec^2(\theta)=\frac{2}{1+u}$$ Invert both sides: $$\displaystyle \cos^2(\theta)=\frac{u+1}{2}$$ Solving for $u$, we find: $$\displaystyle u=2\cos^2(\theta)-1=\cos(2\theta)$$ Hence, we find: $$\displaystyle \tan^{-1}\left(\sqrt{\frac{1-u}{1+u}} \right)=\theta=\frac{1}{2}\cos^{-1}(u)$$ Using the identity $$\displaystyle \sin^{-1}(u)+\cos^{-1}(u)=\frac{\pi}{2}$$ we finally have: $$\displaystyle \tan^{-1}\left(\sqrt{\frac{1-u}{1+u}} \right)=\frac{1}{2}\left(\frac{\pi}{2}-\sin^{-1}(u) \right)$$ Utilizing these results, we now may state: $$\displaystyle \sum_{n=1}^{\infty}\frac{1}{n^2}=4\left(\int_0^{ \frac{1}{2}}\frac{1}{\sqrt{1-u^2}}\sin^{-1}(u)\,du+\frac{1}{2}\int_{ \frac{1}{2}}^1 \frac{1}{\sqrt{1-u^2}}\left(\frac{\pi}{2}-\sin^{-1}(u) \right)\,du \right)$$ Now, let's look at the first integral: $$\displaystyle \int_0^{ \frac{1}{2}}\frac{1}{\sqrt{1-u^2}}\sin^{-1}(u)\,du$$ Using the substitution: $$\displaystyle \alpha=\sin^{-1}(u)\,\therefore\,d\alpha=\frac{1}{\sqrt{1-u^2}}\,du$$ we now have: $$\displaystyle \int_0^{\frac{\pi}{6}}\alpha\,d\alpha=\frac{1}{2} \left[\alpha^2 \right]_0^{\frac{\pi}{6}}= \frac{1}{2}\left(\frac{\pi}{6} \right)^2=\frac{\pi^2}{72}$$ Next, let's break the second integral into two parts: i) $$\displaystyle \frac{\pi}{4}\int_{\frac{1}{2}}^1 \frac{1}{\sqrt{1-u^2}}\,du=\frac{\pi}{4}\left[\sin^{-1}(u) \right]_{\frac{1}{2}}^1=\frac{\pi}{4}\left(\frac{\pi}{2}-\frac{\pi}{6} \right)=\frac{\pi^2}{12}$$ ii) $$\displaystyle -\frac{1}{2}\int_{\frac{1}{2}}^1 \frac{1}{\sqrt{1-u^2}}\sin^{-1}(u)\,du$$ Using the substitution: $$\displaystyle \alpha=\sin^{-1}(u)\,\therefore\,d\alpha=\frac{1}{\sqrt{1-u^2}}\,du$$ we now have: $$\displaystyle -\frac{1}{2}\int_{\frac{\pi}{6}}^{\frac{\pi}{2}} \alpha\,d\alpha=-\frac{1}{4}\left[\alpha^2 \right]_{\frac{\pi}{6}}^{\frac{\pi}{2}}=-\frac{1}{4}\left(\left(\frac{\pi}{2} \right)^2-\left(\frac{\pi}{6} \right)^2 \right)=-\frac{\pi^2}{18}$$ Thus, putting these results together, there results: $$\displaystyle \sum_{n=1}^{\infty}\frac{1}{n^2}=4\left(\frac{\pi^2}{72}+\frac{\pi^2}{12}-\frac{\pi^2}{18} \right)=4\left(\frac{\pi^2}{24} \right)$$ And, we may then state: $$\displaystyle \sum_{n=1}^{\infty}\frac{1}{n^2}=\frac{\pi^2}{6}$$ Shown as desired. So...agentredlum, you now have the floor. 1 person #### agentredlum Math Team Thanx , it would have taken me 2 hours + to write in rhe details using mathjax with my PS3 controller. Are you concerned at all that when $$x = y = 1$$ we get $$\dfrac{1}{1-1}$$ and the Geometric series does not apply? 1 person #### MarkFL Thanx , it would have taken me 2 hours + to write in rhe details using mathjax with my PS3 controller. Are you concerned at all that when $$x = y = 1$$ we get $$\dfrac{1}{1-1}$$ and the Geometric series does not apply? I see $x$ and $y$ as "approaching" 1 from beneath as given by the definition of the rectangle. 1 person Similar Math Discussions Math Forum Date Calculus Math Books Number Theory Number Theory
2020-04-07T11:05:31
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https://math.stackexchange.com/questions/2832794/does-the-following-object-have-such-property
# Does the following object have such property? Let $S_n$ denote the set of bijections on the set $M = \{1, 2, ... , n\}$. Suppose that a set $\Omega \subset S_n$ satisfies the following condition: there is $k \leq n$ such that, for each nonempty subset $A \subset M$ with $|A| \leq k$ and each $a \in A$ there exist exactly $\frac{|\Omega|}{|A|}$ permutations $\pi \in \Omega$ for which $\min_{x \in A} \pi (x) = \pi (a)$. Is it true, that for every nonempty subset $A \subset M$ with $|A| \leq k$ and every $m \in \{1, 2, ... |A|\}$ there exist exactly $\frac{|\Omega|}{|A|}$ permutations $\pi \in \Omega$ such that $\pi (a)$ is the $m$-th largest element of $\pi (A)$? For small $n$-s the statement seems to be true. However, I have no idea how to prove it in general. Any help will be appreciated. • Do you have an example with $\Omega\subsetneq S_n$ for which your “following condition” is true with $k>2$? Certainly it’s true for $\Omega= S_n$ (for any $k$), and regardless of $\Omega$, if you consider subsets $A=\{i,j\}$ of cardinality $2$, the condition “$\sigma(i)$ is $m$th largest in $\sigma(A)$” is the same as “$\sigma(j)$ is smallest” when $m=1$ and “$\sigma(i)$ is smallest” for $m=2$. So it seems like your conjecture could only have a counterexample with $\Omega\subsetneq S_n$ and $k>2$, but I can’t easily think of a proper subset of $S_n$ with your condition ($k>2$). – Steve Kass Jun 30 at 18:03 • I suspect you want $k$ to be a fixed constant, rather than "there exists $k$". – darij grinberg Jun 30 at 18:21 • @SteveKass: $\Omega = A_n$ (the alternating group) and $k \leq n-2$ always works. – darij grinberg Jun 30 at 18:23 The answer is affirmative. Use induction, say, on $m$. By assumption, the statement holds for $m=1$. Denote for $\pi \in S_n$, $A\subset[n]$, $a\in A$ by $\nu_\pi(A,a)$ the rank of $\pi(a)$ among $\pi(x), x\in A$. Then $$N(A,a,m) = \sum_{\pi\in\Omega}\mathbf{1}_{\nu_\pi(A,a)=m}$$ is the quantity we are interested in. If $m = |A|$, then by induction hypothesis $$N(A,a,m) = |\Omega| - \sum_{i=1}^{m-1} N(A,a,i) = |\Omega| - \sum_{i=1}^{m-1} \frac{|\Omega|}{|A|} = \frac{|\Omega|}{|A|}.$$ Otherwise, if $2\le m<|A|$, write by induction hypothesis $$|\Omega| = (|A|-1)\cdot\frac{|\Omega|}{|A|-1} = \sum_{b\in A\setminus \{a\}} N(A\setminus \{b\},a,m-1)\\ = \sum_{b\in A\setminus \{a\}} \sum_{\pi\in\Omega}\mathbf{1}_{\nu_\pi(A\setminus\{b\},a)=m-1} = \sum_{\pi\in\Omega}\sum_{b\in A\setminus \{a\}}\mathbf{1}_{\nu_\pi(A\setminus\{b\},a)=m-1}.$$ Note that if $\nu_\pi(A\setminus\{b\},a)=m-1$, then either $\pi(b)>\pi (a)$ and $\nu_\pi(A,a)=m-1$ or $\pi(b)<\pi(a)$ and $\nu_\pi(A,a)=m$. Therefore, using the induction hypothesis, $$|\Omega| = \sum_{\pi \in |\Omega|} \mathbf{1}_{\nu_\pi(A,a)=m-1} \sum_{b\in A:\pi(b)>\pi(a) }1 + \sum_{\pi \in |\Omega|} \mathbf{1}_{\nu_\pi(A,a)=m } \sum_{b\in A:\pi(b)<\pi(a) }1 \\ = \sum_{\pi \in |\Omega|} \mathbf{1}_{\nu_\pi(A,a)=m-1} (|A|-m + 1) + \sum_{\pi \in |\Omega|} \mathbf{1}_{\nu_\pi(A,a)=m }(m-1)\\ = (|A|-m + 1) N(A,a,m-1) + (m-1) N(A,a,m)\\ = (|A|-m + 1) \frac{|\Omega|}{|A|} + (m-1) N(A,a,m).$$ Hence, $$N(A,a,m) = \frac{|\Omega|}{|A|},$$ as required. • In fact, the cases $m=|A|$ and $m<|A|$ can be combined: I had a mistake in my draft, where the induction was in $k$ and $m$, and $N(A,a,m)$ was written in terms of $N(A,a,m-1)$ and $N(A\setminus\{b\}, a,m)$, so I put this case in there. Now considering $m=|A|$ individually is no longer needed, but let it be. – zhoraster Jul 17 at 5:24
2018-11-16T22:46:26
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https://cs.stackexchange.com/questions/54990/filling-bins-with-pairs-of-balls
Filling bins with pairs of balls A bin is called full if it contains at least $k$ balls. Our goal is to make as many bins as possible full. In the simplest scenario, we are given $n$ balls and may arrange them arbitrarily. In that case, obviously the best we can do is pick $\lfloor n/k \rfloor$ bins arbitrarily and put $k$ balls in each one of them. I am interested in the following scenario: we are given $n$ pairs of balls. We have to put the two balls of each pair in two different bins. Then, an adversary comes and removes one ball from each pair. What can we do to have the maximum possible number of full bins after the removal? A simple strategy is: pick $\lfloor n/(2k-1) \rfloor$ pairs of bins. Fill each bin-pair with $2k-1$ ball-pairs (each bin contains $2k-1$ balls, one ball from each pair). Then, regardless of what our adversary removes, we have in each bin-pair at least one full bin. Do we have a strategy that achieves a larger number of full bins (more than $\lfloor n/(2k-1) \rfloor$)? • I don't believe so Apr 5 '16 at 17:39 • $n$ is given and $k$ is given? $k$ depends on $n$? – Evil Apr 5 '16 at 23:13 • @EvilJS $n$ and $k$ are given, and are independent. Apr 7 '16 at 9:00 • Does the player place all of his $n$ pairs of balls and then the adversary picks $n$ balls?, or does the player place a pair of balls and then the adversary chooses one from that pair and then the player puts the next pair and the adversary picks one and so on until there are no more pairs of balls to place? Apr 10 '16 at 22:20 • @rotia The player places all of his n pairs of balls, and then the adversary picks n balls. Apr 11 '16 at 10:36 TL;DR -- No, there is no better strategy than the simple strategy. Here is the main idea of the proof. When there are not enough balls, there will be a "ball path" from a $k$-full bin to a bin with at most $k-2$ balls. The adversary can pass a ball from that full bin to that less full bin along that path, which can be done repeatedly until the number of $k$-full bins is reduced. Reformulation in Graph theory Suppose we are given a simple finite graph $G(V,E)$ with a function $w: E\to\Bbb Z_{\ge 0}$. We say there are $w(e)$ balls in edge $e$. Let $E_2$ be the (end-marked edge) set $\{(e,v)| e\in E, v\in e\}$. If $d:E_2\to\Bbb Z_{\ge 0}$ satisfies $w(e)=d(e,v_1) + d(e,v_2)$ for every edge $e=\{v_1,v_2\}$, we say that $d$ is $w$-distributing. Any $w$-distributing function $d$ induces a function, which we use the same symbol, $d:V\to\Bbb Z_{\ge 0}$, $d(v) =\sum_{v\in e}d(e,v)$. We say that $d(v)$ balls are in $v$. Given $k\in\Bbb Z_{\gt0}$, let $F_k(d)=\#\{v\in V|d(v)\ge k\}$, the number of $k$-full vertices by $d$. (Erel-Apass Theorem) For any simple finite graph $G(V,E)$ and $w:E\to\Bbb Z_{\ge0}$, we have $\sum_{e\in E}w(e)\geq (2k-1) \min_{w\text{-distributing }d}F_k(d)$ Imagine each vertex is a bin. For each edge $e=\{v_1,v_2\}$, $w(e)$ ball-pairs are put into $v_1$ and $v_2$, each of which getting $w(e)$ balls. Among these $w(e)$ ball-pairs, the adversary may take away $d(e,v_2)$ balls from $v_1$ and $d(e,v_1)$ balls from $v_2$. The end result is the same as if, given all empty bins initially, for each edge $e=\{v_1, v_2\}$, $w(e)$ balls are put into it and, then, $d(e,v_1)$ and $d(e,v_2)$ balls are distributed to $v_1$ and $v_2$ respectively by the adversary. Hence, Erel-Apass theorem says that in order to ensure $t$ k-full bins after a smart adversary's removal, at least $(2k-1)t$ pairs of balls are needed. In another word, an optimal strategy to have the maximum possible number of full bins left is indeed the "simple strategy", which repeatedly fills a different pair of bins with $2k-1$ ball-pairs until we do not have enough balls to repeat. Proof of the theorem For the sake of contradiction, let $G(V,E)$ and $w$ be a counterexample whose number of vertices is the smallest among all counterexamples. That is, there is $w$-distributing $m$ such that $F_k(m)$ is minimal among all $F_k(d)$ of $w$-distributing function $d$. Furthermore, $$\sum_{e\in E}w(e)\lt (2k-1)F_k(m)$$ Let $V_s=\{v\in V | m(v)\le k-2\}$. Let $V_\ell=\{v\in V|m(v)\geq k\}$. So $F_k(m)=\#V_\ell$. Claim one: $V_s\neq\emptyset$. Proof of claim one. Suppose otherwise that $V_s$ is empty. $$\sum_{v\in V}m(v)= (k-1)\#V +\sum_{v\in V}(m(v)-(k-1)) \ge (k-1)\#V + \#V_\ell \gt (k-1)\#V$$ Let us also reuse $w$ as a function from $V$ to $\Bbb Z_{\ge 0}$ such that $w(v)=\sum_{v\in e}w(e)$ for any $v\in V$. \begin{align} \sum_{v\in V}w(v) &=\sum_{v\in V}\sum_{v\in e }w(e) =\sum_{e\in E }\sum_{v\in e}w(e) =\sum_{e\in E}2w(e) =2\sum_{e\in E}w(e)\\ &=2\sum_{e\in E }\sum_{v\in e}m(e,v) =2\sum_{v\in V}\sum_{v\in e }m(e,v) =2\sum_{v\in V}m(v)\\ &\gt 2(k-1)\#V \end{align} So there must be a vertex $b$ such that $w(b)\ge 2k-1$. Consider the induced setup $G'(V', E')$ and $w'$, where $V'=V\setminus\{b\}$, $G'(V', E')$ is the induced graph $G[V']$ and where $w'=w|_{E'}$. For any $w'$-distributing function $d'$, we can extend it to a $w$-distributing function $d_{d'}$ where $d_{d'}$ is the same as $d'$ on $E'$ while $d_{d'}(e, b)=w(e)$ for every edge $e$ adjacent to $b$. Note that $F_k(d_{d'})= F_k(d') + 1$ since $d_{d'}(b)=\sum_{b\in e}d_{d'}(e,b) = \sum_{b\in e}w(e)=w(b)\ge 2k-1\ge k$. Then \begin{align}\sum_{e\in E'}w'(e)&\le\sum_{e\in E}w(e) - w(b)\\ &\lt (2k-1)F_k(m) -(2k-1)\\ &=(2k-1) \left(\min_{w\text{-distributing }d}F_k(d) -1\right)\\ &\le (2k-1) \left(\min_{w'\text{-distributing }d'}F_k(d_{d'})-1\right)\\ &\le (2k-1) \min_{w'\text{-distributing }d'}F_k(d') \end{align} So, $G'(V', E')$ and $w'$ is a counterexample whose number of vertices is smaller than the number of vertices in $G$. That cannot true by our assumption about $G(V,E)$ and $w$. So claim one is proved. For any vertex $v$, define $v$ $d$-reachable from vertex $u$ if there is a path $u_0= u, u_1, u_2, \cdots, u_m, u_{m+1}=v$, $m\ge0$ such that $d(\{u_i, u_{i+1}\}, u_i)>0$. Let $V_r=V_\ell\cup \{v\in V|\exists u\in V_\ell \text{ and } v\text{ is } m\text{-reachable from } u \}$. Claim two: $V_r = V$ Proof of claim two: Suppose $V_r\neq V$. For any vertex $v\in V_r$ and $u\notin V_r$, since we cannot reach $u$ from $v$, if $\{v, u\}$ is an edge, then $w(\{v, u\}, v) = 0.$ Consider the induced setup $G'(V', E')$ and $w'$, where $v'=V_r$, $G'(V', E')$ is the induced graph $G[V']$ and where $w'=w|_{E'}$. For any $w'$-distributing function $d'$, we can extend it to a $w$-distributing function $d_{d'}$ where $d_{d'}$ is the same as $d'$ on $E'$ and the same as $m$ on other edges. Note that $F_k(d_{d'})= F_k(d')$ since all vertices with no less than $k$ balls inside are in $V_\ell\subset V_r$. Then \begin{align}\sum_{e\in E'}w'(e)&\le\sum_{e\in E}w(e)\\ &\lt (2k-1)F_k(m)\\ &=(2k-1) \min_{w\text{-distributing }d}F_k(d)\\ &\le (2k-1) \min_{w'\text{-distributing }d'}F_k(d_{d'})\\ &\le (2k-1) \min_{w'\text{-distributing }d'}F_k(d') \end{align} So, $G'(V', E')$ and $w'$ would be a counterexample whose number of vertices is smaller than the number of vertices in $G$. That cannot be true by our assumption about $G(V,E)$ and $w$. So claim two is proved. Now let us prove the theorem. Since $V_r=V$ and $V_s\neq\emptyset$, there is a path $u_0= u, u_1, u_2, \cdots, u_m, u_{m+1}=v$, $m\ge0$ with $m(u)\gt k$, $m(v)\leq k-2$ and $d(\{u_i, u_{i+1}\}, u_i)>0$. Let us construct a new $w$-distributing function $r(m)$ from $m$ so that $$r(m)(e, u)= \begin{cases} m(\{u_i, u_{i+1}\}, u_i) -1 & \text{ if } (e,u)=(\{u_i, u_{i+1}\},u_i)\text { for some } 0\le i\le m\\ m(\{u_i, u_{i+1}\}, u_{i+1}) +1 & \text{ if } (e,u)=(\{u_i, u_{i+1}\},u_{i+1})\text { for some } 0\le i\le m\\ m(e,u) &\text{ otherwise } \end{cases}$$ $m$ and $r(m)$ agrees on all vertices except $v$ and $u$, $m(v)\lt r(m)(v)\le k-1$ and $r(m)(u)\lt m(u)$. We can apply this procedure on $r(m)$ to get $r^2(m)$. Repeating this $i$ time for some large enough $i$, we will obtain a $w$-distributing function $r^i(m)$ with $F_k(r^i(m))=0$. However, we have assumed that $F_k(m)>0$ is the minimum among $F(d)$ of $w$-distributing function $d$. This contradiction shows that we have proved the Erel-Apass theorem. • I read the proof, it looks good. In fact, if I understand correctly, it is even more general since it allows for an arbitrary graph - my question is a special case where G is the complete graph. Is this correct? Another question: where exactly does the proof use the fact that m is such that Fk(m) is minimal? I see that it is used only at the last paragraph - are the previous claims in the proof true without this fact? Oct 31 '18 at 8:03 • Yes, the theorem is correct for any graph since it says "for any (simple finite) graph G(V,E)". The minimality of $F_k(m)$ is necessary for each claim. If you search for "counterexample", you will find where the minimality is used. Nov 5 '18 at 6:17
2022-01-29T05:24:52
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https://math.stackexchange.com/questions/2918124/is-empty-set-with-usual-metric-space-x-mathbbr-d-bounded
# Is empty set with usual metric space $(X=\mathbb{R},d)$ bounded? Here R is set of real numbers and $d$ is distance metric i.e $|x-y|$. I wanted to find if empty set under given metric space is compact and we know by Heine Borel theorem that a set is compact iff it is closed and bounded. Empty set is closed because its complementary set i.e. set of real numbers $\mathbb{R}$ is open.Hence i just need to prove if the empty set is bounded. The problem i am facing is , how can we define $|x-y|$ on empty set , to find $\exists m$ , $m \in \mathbb{R}$ , such that $|x-y|\leqslant m$. • Write out the definition carefully with quantifiers. – Ted Shifrin Sep 15 '18 at 18:13 • A set $E$ is bounded iff there exists $m>0$ such that, for all $x,y \in E$ we have $d(x,y) \le m$. So what do you get when $E = \varnothing$? – GEdgar Sep 15 '18 at 18:15 • It is vacuously bounded. – Eduardo Longa Sep 15 '18 at 18:15 • @GEdgar I cannot have any pair $x,y$ such that$d(x,y)<=m$ , hence i cannot say it is bounded or not bounded. – user585269 Sep 15 '18 at 18:20 • Let $E = \emptyset$, and $m > 0$. Can you find a pair $x,y \in E$ such that $d(x,y) > m$? – Theoretical Economist Sep 15 '18 at 18:23 This is answered in comments already, but here's an honest-to-god StackExchange answer. Your question can be asked more succinctly as: "Is $\{\}$ bounded when considered as a subset of $\mathbb{R}$ with its standard metric?". The answer is Yes. In fact, the empty set is bounded when thought of as a subset of any metric space. To see why, you have to write down a definition of boundedness. It should probably come out as this: $X$ is bounded with respect to the metric $d$ if there exists a number $M \in \mathbb{R}$ such that, $d(a,b) < M$ for all $a, b \in X$. Side note: I would have got equivalent definitions if I'd said $M \in \mathbb{N}$ or asked that $d(a, b) \le M$. You can check you understand the definition above by checking you understand why these definitions would give the same thing. I claim that $\{\}$ is bounded with respect to any $d$. To prove this, I can set $M$ to be any number. After all, I just need to prove that if you choose $a, b \in \{\}$ then $d(a,b) < M$. But you can't choose find such $a$ or $b$ because the empty set is... empty! As such, the inequality is never actually checked. As you correctly wrote in the question, you can now use the Heine-Borel theorem to show that the empty set is compact as a subset of $\mathbb{R}$. In fact, the Heine-Borel theorem is quite heavy machinery to use to prove the compactness of the empty set. If you think for a minute about the definition of compactness based on finite subsets of open covers, you should be able to prove directly that the empty set is compact as a subset of any topological space. • To the proposer: If $S$ is any statement, phrase, clause, formula, etc., and if $X=\emptyset$ then the sentence $\forall a\in X\;(S)$ is true. Because its negation is $\exists a\in X\;(\neg S),$ which cannot be true unless some $a$ exists that belongs to $X.$ – DanielWainfleet Sep 16 '18 at 1:01
2019-08-25T00:56:44
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https://math.stackexchange.com/questions/2430323/what-is-the-likelihood-of-two-line-segments-crossing
# What is the likelihood of two line segments crossing? Consider a square space. Randomly select 4 points. Randomly connect two sets of two to each other with line segments. What is the chance of the line segments intersecting? (I will maybe try and solve this with a little Monte-Carlo simulation but would be very interested in an analytic solution.) • Do you mean lines, or line segments (i.e. are the lines infinitely long in both directions, or do they stop at the points)? – Arthur Sep 15 '17 at 8:27 • The title was unclear. You are right. I updated it – Christian Sep 15 '17 at 8:31 • I went through some coding. I may be incorrect but the estimated probability is around $0.111$... May anyone confirm this result? – Wyllich Sep 15 '17 at 9:39 • @Wyllich I may be wrong, but it seems to me that the analytical solution (see answer below) yields a result of 25/108, which is about twice the probability you obtained. – Anonymous Sep 15 '17 at 10:56 • @Anonymous Yes. You are right. I checked my code and indeed, I was badly checking if the calculated intersection point belongs to both lines or not. I agree with $\approx 0.23$ too. – Wyllich Sep 15 '17 at 11:29 Let's be perfectly clear about the assumptions. I assume that "randomly selecting 4 points" A,B,C,D in the unit square means selecting uniformly and independently at random their coordinates $x_A,y_A,x_B...$ in the $[0,1]$ range; and that one of the 6 possible pairings between points to form the segments is also chosen uniformly and independently at random -- so the problem is equivalent to asking what is the probability that, e.g. AC and BD intersect. It's easy to see that the solution is $\frac{25}{108}$, or about $23\%$, if one is willing to take as a given the solution for Sylvester's four-point problem for the square -- the probability that the convex hull of four points chosen uniformly and independently at random in a square is a quadrilateral, which is $\frac{25}{36}$. It is immediate to see that two segments chosen by a random pairing of four points intersect if, and only if, the convex hull of the four points is indeed a quadrilateral, and the two segments are its diagonals. If the former condition is satisfied, the latter happens with probability $\frac{1}{3}$ (the probability that a given point is paired with "middle one" of the remaining three), yielding a solution of $\frac{25}{36}\cdot\frac{1}{3}=\frac{25}{108}$ for the original problem. Sylvester's four-point problems allows one to obtain the answer in exactly the same way even if the points are chosen uniformly at random in another convex region, simply by multiplication by $\frac{1}{3}$. The general formula for Sylvester's four-point problem in a convex region $R$, as given in the above link, equals $1-4\frac{\bar{A_R}}{A(R)}$ where $A(R)$ is the area of the region (in the case of the unit square, $1$) and $\bar{A}_R$ equals the expected area of a triangle obtained from $3$ points chosen uniformly at random within it. The formula is easily obtained noting that probability that one given point lies within the triangle formed by the other $3$ is $\frac{\bar{A_R}}{A(R)}$, and since the events of this happening for different points are mutually exclusive, the probability that no point lies within the triangle formed by the other three is indeed $1-4\frac{\bar{A}_R}{A(R)}$. Computing $\frac{\bar{A_R}}{A(R)}$ is non-trivial, in general. If R is a square (our case), a simple proof that $\frac{\bar{A_R}}{A(R)}=\frac{11}{144}$ can be found here. More in general if R is regular polygon of $n$ sides the solution is given by Alikosky's formula, $\frac{\bar{A_R}}{A(R)}=\frac{9\cos^2 (2\pi/n)+52\cos(2\pi/n)+44}{36n^2\sin^2(2\pi/n)}$. Note that this is a strictly decreasing function of $n$, so the intersection probability is strictly increasing with $n$, and the limit for $n\to\infty$, i.e. $\frac{35}{48\pi^2}$, yields the solution when picking points in the circle. So, in general, the probability of intersection of $2$ segments, each made by $2$ points chosen uniformly and independently at random from a regular $n$-agon or any of its affine transformations (such as a rectangle or parallelogram) is: $\frac{1}{3}\left(1-4\cdot\frac{9\cos^2(2\pi/n) + 52\cos(2\pi/n) + 44}{36 n^2\sin^2 (2\pi/n)}\right)$ with the limit for $n\to\infty$, i.e. $\frac{1}{3}-\frac{35}{36\pi^2}$, yielding the probability when choosing the points from a circle (or any affine transformation, such as an ellipse). • You mention the probability of the square being quadrilateral as $\frac{25}{36}$. Is this not the probability of it being convex? I think it will always be a quadrilateral according to the accepted definition. mathworld.wolfram.com/Quadrilateral.html Or am I misreading something? – Christian Sep 25 '17 at 11:11 • No, I am saying that 25/36 is the probability that the convex hull of four points, chosen uniformly at random within the square, is a quadrilateral (and thus a convex quadrilateral). Note that with some probability (11/36) the convex hull is a triangle instead, i.e. one of the four points falls inside the triangle formed by the other three. – Anonymous Sep 25 '17 at 23:26 First let's select 4 points inside a square. Let $P_c$ be a probability that these 4 points form a convex figure (none of the selected points is inside the triangle formed by the other three). If this is the case, the probability that two line segments intersect is 1/3. Otherwise the probability is 0. So, the probability we are looking for is: $P=1/3 * P_c$ Now we need to find the $P_c$. Let's calculate the probability $P_1$ that the first point is inside the triangle formed by the three others. It is $P_1 = S_{234}/S_{square}$, where $S_{234}$ is the expected area of the triangle formed by points 2, 3 and 4. Same for $P_2$, $P_3$ and $P_4$. Obviously $P_1=P_2=P_3=P_4$. Probability that any point is inside of triangle formed by other 3 is $P_1+P_2+P_3+P_4$ because no more than 1 of these outcomes can happened simultaneously. $P_c = 1 - 4*P_1 = 1 - 4 * S_{234}$ And now we only need to find $S_{234}$ - expected area of a triangle formed by randomly selected 3 points inside a square. The expected area of a triangle formed by three points randomly chosen from the unit square They say it is $11/144$. $P=1/3 * (1 - 4*11/144) = 25/108$ This is in a good agreement with my experiments. This is my approach: Consider a line $y=mx+c$ that divides the unit square into two trapeziums. Let us call the left trapezoid $a$, and the right trapezoid $b$. By solving the equations $(y=mx+c)=1$ and $(y=mx+c)=0$, or $mx+c=1$ and $mx+c=0$, the $x$-values of the intersections are $x = \frac{-c+1}{m}$, and $\frac{-c}{m}$. Since the $y$-values are $0$ and $1$, the intersection coordinates are $(\frac{-c}{m}, 0)$ and $(\frac{-c+1}{m}, 1)$. Using this information, the area of trapezium $a$ is $1 * \frac{(-c/m+1)+(-c/m)}{2}$, and of $b$ is $1 - (1 * \frac{(-c/m+1)+(-c/m)}{2})$ (since $a+b=1$). Simplifying gives the area of $a$ as $-\frac{c}{m}+ \frac{1}{2}$, and the area of $b$ as $\frac{c}{m}+\frac{1}{2}$. For the line connecting two points to not cross the green line, the points cannot be in $a$ and $b$. Since using geometric probability, the probability of this is $ab-ba$, we can subtract this condition from $1$ to get $1 - 2(-\frac{c}{m}+ \frac{1}{2})(\frac{c}{m}+\frac{1}{2})$, and after simplifying we have the probability as $1 - 2(-\frac{c^2}{m^2}+\frac{1}{4})$ or $$\frac{2c^2}{m^2}+\frac{1}{2},$$ where $c$ is the $y$-intercept of the line (when extended), and $m$ is its slope. Note: This example doesn't work when $c=0$, because regardless of what $m$ is chosen, it would equal $0+\frac{1}{2}=\frac{1}{2}$, and when the line crosses a point not between $0$ and $1$ when $x=0$ or $x=1$. • Sorry to all of you reading my answer - I just realised the final answer doesn't make any sense. I will try to correct the error within $10$ minutes as much as I can. – Toby Mak Sep 15 '17 at 8:58 • The question asks for line segments, not lines that intersect inside of the square – Dominik Sep 15 '17 at 9:01 • Lines and line segments are the same things, except that one is a shortened version of another. I made sure in my answer that I did not include any solutions with intersections outside the square, for example by setting $a = 1-b$. It doesn't matter whether they are lines or line segments - their equations are essentially very similar. – Toby Mak Sep 15 '17 at 9:08 • Consider the line segments through $(0, 0.5)$, $(0.5,0.5)$ resp. $(0.75, 0.25)$, $(0.75, 0.75)$. They do not intersect, but would intersect according to yourcalculations. – Dominik Sep 15 '17 at 9:11 • That's true, but your second line is vertical, which cannot be expressed in any equation because of division by zero. You can't prove it's just my formula that doesn't work. – Toby Mak Sep 15 '17 at 9:15 If in four points if $$\color{red}{1]}$$ all are collinear, line segments will never 'intersect' $$\color{red}{2]}$$ Three are collinear, still no intersection. $$\color{red}{3]}$$ Two points are always collinear:) $$..................................$$ 4 non collinear points will form quadrilateral. $$\color{red}{\text{if quadrilateral is convex}}$$ Lines segments drawn must be either sides or diagonals. Only diagonals intersect. $$\color{red}{\text{if quadrilateral is concave}}$$ No intersection Solve using this information • Collinear happens with probability zero. Also, I think the odd formatting makes the post much more difficult to read. – Hurkyl Sep 15 '17 at 9:33 • What is the probability of convex quadrilateral? – neonpokharkar Sep 15 '17 at 9:49
2019-08-25T19:56:04
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https://math.stackexchange.com/questions/2526636/proving-lim-limits-n-to-infty-frac8n2-54n27-2-using-epsilon-d
# Proving $\lim\limits_{n\to\infty} \frac{8n^2-5}{4n^2+7} = 2$ using $\epsilon-\delta-$ definition. I'm trying to prove the limit of this sequence using the formal definition. I've looked at other questions on the site but from the ones I've seen, the $n^2$ term always seems to cancel out, making it simpler. Show that $$\lim_{n\to\infty} \frac{8n^2-5}{4n^2+7} = 2$$ So this is how I started: $$\left|\frac{8n^2-5}{4n^2+7} -2 \right| = \left|\frac{-19}{4n^2+7}\right| = \frac{19}{4n^2+7} \leq \frac{19}{4n^2} = \epsilon$$ Let $\epsilon > 0 \implies n = \sqrt{\frac{19}{4\epsilon}}$ By Archimedian Property: $N > \sqrt{\frac{19}{4\epsilon}}$ If $n \geq N \geq \sqrt{\frac{19}{4\epsilon}}$ $$\left|\frac{8n^2-5}{4n^2+7} -2 \right| \leq \frac{19}{4n^2} < \frac{19}{4(\frac{19}{4\epsilon})} = \epsilon$$ • Looks correct to me. – user491874 Nov 18 '17 at 21:23 • Welcome to math.stackexhange! This looks fine. – Hans Engler Nov 18 '17 at 21:24 • +1 for using MathJax and showing your attempted solution. – Remy Nov 18 '17 at 21:26 • Yes, it’s correct. Just one small point, use “there exist” before N in the step where you use Archimedian property. – Mayuresh L Nov 18 '17 at 21:28 • Canceling out the $n^2$ will leave $\frac 5{n^2}$ and $\frac 7{n^2}$. I wouldn't say that is any easier. Maybe more intuitive where to go, but not easier.... any way, what you did was exactly what you were supposed to do and you did it exactly right. (Ithink... $\frac {8n^2 - 5}{4n^2 + 7} - 2 = \frac {(8n^2 - 5) - (8n^2 + 14)}{4n^2 + 7} = \frac {-19}{4n^2 + 7} ....$ Yep, it checks out. (Lack of an $n$ term threw me but it's fine...) – fleablood Nov 18 '17 at 22:35 we have $$\left|\frac{8n^2-5}{4n^2+7} -2 \right| = \left|\frac{-19}{4n^2+7}\right| < \frac{19}{4n^2} =$$ Now let such that $$\frac{19}{4n^2} < \epsilon\Longleftrightarrow n>\sqrt{\frac{19}{4\epsilon}}$$ Now choosing, $$N =\left\lfloor \sqrt{\frac{19}{4\epsilon}}\right\rfloor +1$$ then we have $$n > N=\left\lfloor \sqrt{\frac{19}{4\epsilon}}\right\rfloor +1\Longrightarrow n>\sqrt{\frac{19}{4\epsilon}}\\\Longrightarrow \epsilon>\frac{19}{4n^2} > \left|\frac{8n^2-5}{4n^2+7} -2 \right|$$ that is we have $$n > N=\left\lfloor \sqrt{\frac{19}{4\epsilon}}\right\rfloor +1\Longrightarrow \left|\frac{8n^2-5}{4n^2+7} -2 \right|< \epsilon$$
2021-05-08T20:30:09
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https://math.stackexchange.com/questions/3839244/is-3-a-subset-of-1-1-2-1-2-3/3839249
# Is $\{3\}$ a subset of $\{\{1\},\{1,2\},\{1,2,3\}\}$? Is $$\{3\}$$ a subset of $$\{\{1\},\{1,2\},\{1,2,3\}\}$$? If the set contained $$\{3\}$$ plain and simply I would know but does the element $$\{1,2,3\}$$ include $$\{3\}$$ such that it would be a subset? • No, it's not. Let $A=\{\{1\},\{1,2\},\{1,2,3\}\}$. The subsets must themselves contain sets, since A is a set of sets. Then $\{\{1\}\}$ is a subset of A, $\{\{1\},\{1,2\}\}$ is another subset of A. $\{\{1,2\},\{1,2,3\}\}$ is also a subset of A. $\{3\}$ is a subset of an element of A. Sep 24 '20 at 21:03 • No, when you have a set of sets, you do not "look through" the brackets. The sets are members, but the members of those sets are not members. Sep 24 '20 at 21:04 • Similar question Sep 24 '20 at 21:05 • "Is $3=\{1\}$?" ; "Is $3=\{1,2\}$?" ; "Is $3=\{1,2,3\}$?": If the answer to one of these questions is yes, then $\{3\}$ is a subset of $\{\{1\},\{1,2\},\{1,2,3\}\}$; otherwise, it isn't. – user239203 Sep 24 '20 at 21:05 • @Gae.S.And it should be noted that the answer may be yes! Sep 24 '20 at 21:13 Let $$X$$ be a set. We say $$Y \subseteq X$$ ($$Y$$ is a subset of $$X$$) if, for all $$x \in Y$$, we have $$x \in X$$. Examine the sets $$Y = \{3\}$$, $$X = \{\{1\},\{1,2\},\{1,2,3\}\}$$. Take $$x = 3 \in Y$$. Is $$3 \in X$$? Trickier problem: If $$X = \{\{1\},\{3\},\{1,2\},\{1,2,3\}\}$$, is $$3 \in X$$? • Why is the trickier problem trickier? $3\not \in X$ because $3\ne \{1\};3\ne \{3\},3\ne \{1,2\}, 3\ne \{1,2,3\}$. Why is that trickier. Sep 24 '20 at 21:21 • I would expect someone to say "Oh, I see 3 by itself" and say it is in the set. I wanted to emphasize the difference between $\{3\}$ and $3$. Sep 24 '20 at 21:23 No. the elements inside elements of the set do not count. The elments of your big set are: • $$\{1\}$$ • $$\{1,2\}$$ • $$\{1,2,3\}$$. The elements of your small set are: • $$3$$ So $$\{3\}$$ is a subset only if $$3$$ is equal (the same thing; !!!!!NOT!!!! an element within) one of the elements $$\{1\}$$ or $$\{1,2\}$$, or $$\{1,2,3\}$$. But none of those are the same thing as $$3$$ so $$\{3\}$$ is not a subset. But in some text the natural numbers are defined as $$0 = \emptyset$$ $$1= \{\emptyset\}$$ $$2= \{\emptyset, 1\}$$. $$3 = \{\emptyset, 1, 2\}$$ So we could have a trick thing of $$\{3\} \subset \{\{\emptyset,1\}, \{\emptyset,1,2\}, \{\emptyset,1,2,3\}\}$$ not because $$3 \in \{\emptyset,1,2,3\}$$ (that's utterly irrelevent), but because $$3 = \{\emptyset, 1, 2\}$$ and the set $$\{\{\emptyset,1\}, \{\emptyset,1,2\}, \{\emptyset,1,2,3\}\}$$ is equally equal to the set $$\{2,3,4\}$$ and $$\{3\}\subset \{2,3,4\}$$. No it is not, we have that $$\{1\}\in \{\{1\},\{1,2\},\{1,2,3\}\}$$ but $$\{3\}$$ is not an element of the set, what is true is that $$3\in\{1,2,3\}$$. As you can see, there are plenty of good answers here. I think that the most important thing to understand is: $$3 \neq \{ 3 \}.$$ These are two different objects. Moreover the followings hold: $$3 \in \{3\}$$ $$3 \in \{\ldots, 3, \ldots \}$$ $$\{3\} \in \{\ldots, \{3\}, \ldots\}$$ and $$\{3\} \not\in \{3\}$$ $$\{3\} \not\in \{\ldots, 3, \ldots \}$$ • In for a penny..... might as well add (and perhaps germaine to the common misconception) $3\not \in \{........., \{3\},......\}$ (unless it is one of the objects omitted in the "..."s.) Sep 24 '20 at 21:59
2021-10-27T06:26:43
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http://mathoverflow.net/questions/75447/computing-permutations-with-partial-duplicates?sort=votes
# Computing Permutations with Partial Duplicates I am looking for a way to compute the number of $K$ permutations of a multiset with $N*D$ elements where each group has exactly $D$ equal elements (and typically $D < N$ ). I've got an application that actually generates these unique permutations and works on them, but I'd like to understand how I can compute the number of sets I'll have across various inputs without computing the entire result. Example (in R): N <- 19 K <- 4 # Implied D = 3 by just duplicating it in-place three times. a <- append(1:N, append(1:N, 1:N)) b <- unique(gtools::permutations(length(a), K, a, set=FALSE)) nrow(b) in this case will be 130,302. This is slow and inelegant. Can someone help me do this with actual math? Expanding a bit If N is 9 and D is 3, my input might look like this: 1 1 1 2 2 2 3 3 3 4 4 4 5 5 5 6 6 6 7 7 7 8 8 8 9 9 9 A standard permutation would look like this: 1 1 1 2 1 1 1 2 1 1 1 2 1 1 1 3 1 1 1 3 1 1 1 3 But at this point, I want to treat the things that look the same as the same, so I deduplicate to get the following: 1 1 1 2 1 1 1 3 1 1 1 4 1 1 1 5 1 1 1 6 1 1 1 7 The first (full permutation) provides 421,200 rows: $(9*3)! \over (9 * 3 - 4)!$ My final, deduplicated answer is 6,552 rows. I'd like to know how I can get that without generating them all. New Discovery For my initial case where $D = K - 1$, I get the correct answer with $N^K - N$. - what's a "unique set of k "? could you make an example? –  Pietro Majer Sep 14 '11 at 21:18 does this help? –  dustin Sep 14 '11 at 21:31 It looks like you want to count ordered $k$-tuples from a set of size $n$ with at most $d$ repetitions of each element. –  Douglas Zare Sep 15 '11 at 0:30 The ordering part doesn't matter much to me (depending on what you meant by it). I want AAAB and ZZZY and every permutation between except where that permutation is identical to another. –  dustin Sep 15 '11 at 3:49 Do you want AAAB to be different from AABA or not? –  Douglas Zare Sep 15 '11 at 4:22 You could try to use exponential generating functions. For each of the N letters you could use the exponential generating function $$\sum_{i=0}^D \frac{x^i}{i!}$$ Cause each letter can be used at most D times this is the same for each letter. Then for using all different letters the egf's have to be multipled (you have N different letters so N times): $$\left(\sum_{i=0}^D \frac{x^i}{i!}\right)^N$$ Then you are looking for the amount of different words of length K which is if you expand the expression above (which is possible if you put into it values for $D, N$. Then the coefficient of $x^K$ multiplied by $k!$ is your solution. What I wasn't able to do right now is trying to expand the generating function into a series without using concrete values for $D$ and $N$. - I'm thinking the result should be: n ^ k for d >= k n ^ k - n ^ (k - d) for d < k - It's not deduplicated in the sense that the same item doesn't show up, but that the same sequence doesn't show up in the result. i.e. AAAB, AAAC, but not another AAAB where the As were rearranged. –  dustin Sep 14 '11 at 23:10 I don't know why I was making this far too hard. D will not matter, because no matter what it is, it will get deduplicated. Since what you're doing is (n*d)Pk with deduplication, your result is just n^k. –  nullghost Sep 15 '11 at 0:00 Okay, I missed the bit where you can only choose a number a maximum of D times. If D is greater than or equal to K, number of rows = N^K If D is less than K, number of rows should = N^K - N^(K-D) Test that for me, see if the results make sense. –  nullghost Sep 15 '11 at 0:08 That's off in just a couple of cases: pastebin.com/LUd8BMSR –  dustin Sep 15 '11 at 1:03 The number of $K$-permutations of the numtiset $\{ 1^D, 2^D, \ldots, N^D \}$ is $$\sum_{j_1+j+2+\dots+j_N=K\atop 0 \leq j_i \leq D} \binom{K}{j_1,j_2,\dots,j_N}.$$ (summands here are multinomial coefficients) Alternatively, denoting by $m_i$ the number of $j$'s equal $i$, we get a formula as the sum over restricted partitions: $$\sum_{0m_0+1m_1 + \dots + Dm_D = K\atop m_0 + m_1 + \dots + m_D = N,\quad m_i\geq 0} \binom{N}{m_0,\dots,m_D} \frac{K!}{1!^{m_1} 2!^{m_2} \cdots D!^{m_D}}$$ For the example with $N=9$, $D=3$, $K=4$, the latter formula consists of four summands and gives: $$\binom{9}{5,4} \frac{4!}{1!^4} + \binom{9}{6,2,1}\frac{4!}{1!^2 2!^1} + \binom{9}{7,1,1}\frac{4!}{1!^1 3!^1} + \binom{9}{7,2}\frac{4!}{2!^2}$$ $$= 3024 + 3024 + 288 + 216 = 6552$$ as expected. -
2014-10-23T08:26:04
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https://mathematica.stackexchange.com/questions/237692/simplifying-code-that-displays-a-square-rolling-inside-a-circle
# Simplifying code that displays a square rolling inside a circle The vertex A,B of the square ABCD with side length $$\sqrt{2}$$ is on the circle with radius $$\sqrt{2}$$, vertex C,D is inside the circle, roll the square ABCD along the inside of the circle counterclockwise without sliding. My code is already working, but I believe he can be more concise Clear["*"]; p0 = N @ {{-√2/2, -√6/2}, {√2/2, -√6/2}, {√2/2, √2-√6/2}, {-√2/2, √2-√6/2}}; p1[t_] := RotationTransform[t,p0[[1]]][p0]; p2[t_] := RotationTransform[t,p1[π/6][[4]]][p1[π/6]]; p3[t_] := RotationTransform[t,p2[π/6][[3]]][p2[π/6]]; p4[t_] := RotationTransform[t,p3[π/6][[2]]][p3[π/6]]; p5[t_] := RotationTransform[t,p4[π/6][[1]]][p4[π/6]]; p6[t_] := RotationTransform[t,p5[π/6][[4]]][p5[π/6]]; Manipulate[ Graphics[ {Circle[{0, 0}, Sqrt[2]], EdgeForm[Black], Opacity[0.1], Polygon[ Which[ t < π/6, p1[t], t < 2π/6, p2[t-π/6], t < 3π/6, p3[t-2π/6], t < 4π/6, p4[t-3π/6], t < 5π/6, p5[t-4π/6], True, p6[t-5π/6]]]}, PlotRange -> 2], {t, 0, 6π/6}] Some attempts, but lacks the rotation process With[{p0=N@{{-√2/2,-√6/2},{√2/2,-√6/2},{√2/2,√2-√6/2},{-√2/2,√2-√6/2}}}, Manipulate[With[{pts=Fold[RotationTransform[π/6,#[[-Mod[#2,4,1]]]]@#&,p0,Range[0,i]]}, Graphics[{Circle[{0,0},Sqrt[2]],{EdgeForm[Black],Opacity[0.1],Polygon@pts}}, PlotRange->2]],{i,0,5}]] Do you have a better way? Tips that may be useful Which[t < π/6, t, t < (2 π)/6, t - π/6, t < (3 π)/6, t - (2 π)/6, t < (4 π)/6, t - (3 π)/6, t < (5 π)/6, t - (4 π)/6, True, t - (5 π)/6] == Mod[t, Pi/6] • Just transform a Rectangle---not each point that defines one. Jan 5, 2021 at 18:48 EDIT 1. Fixed kernel crashing. The problem was WhenEvent choking, because it was continually checking whether the current pivot (which lies exactly on the bounding circle) goes outside the bounding circle. EDIT 2. Here is a port to GNU Octave: https://github.com/yawnoc/tumbling-polygon # System of ODEs in the complex plane For a completely different take, observe that rigid rotation of the $$n$$th vertex $$z_n$$ around a pivot $$p$$ can be expressed as the differential equation $$\frac{\mathrm{d}z_n}{\mathrm{d}t} = \mathrm{i} \,(z_n - p),$$ where $$\mathrm{i} = \sqrt{-1}$$, and we have scaled out the angular velocity. This is because the instantaneous velocity is (1) at right angles to the displacement vector $$z_n - p$$, and (2) proportional to the distance $$|z_n - p|$$. (Remember multiplication by $$\mathrm{i}$$ effects a 90-degree rotation.) If the vertex $$z_n$$ happens to coincide with the pivot location $$p$$, this nicely reduces to $$\mathrm{d}z_n / \mathrm{d}t = 0$$. So we solve a system of such ODEs in $$(z_1, z_2, z_3, z_4, p)$$, where we handle the collision with the bounding circle $$|z| = \sqrt{2}$$ using WhenEvent; if $$|z_n| > \sqrt{2}$$, we set $$p$$ to $$z_n$$. ClearAll["Global*"]; vertexInitialPositions = Complex @@@ N @ { {-Sqrt[2]/2, -Sqrt[6]/2}, {Sqrt[2]/2, -Sqrt[6]/2}, {Sqrt[2]/2, Sqrt[2] - Sqrt[6]/2}, {-Sqrt[2]/2, Sqrt[2] - Sqrt[6]/2}, Nothing }; vertexCount = Length[vertexInitialPositions]; vertexRotationEquations = Table[z[i]'[t] == I (z[i][t] - p[t]), {i, vertexCount}]; pivotFixtureEquation = p'[t] == 0; vertexInitialConditions = Table[z[i][0] == vertexInitialPositions[[i]], {i, vertexCount}]; pivotInitialCondition = p[0] == First[vertexInitialPositions]; boundingRadius = Sqrt[2]; collisionHandling = Table[ whenEventDummy[ (* If this vertex isn't pivot AND goes outside bounding circle *) z[i][t] != p[t] && Abs @ z[i][t] > boundingRadius, (* Set the pivot to this vertex *) p[t] -> z[i][t] ] , {i, vertexCount} ] /. {whenEventDummy -> WhenEvent}; dependentVariables = Table[z[i], {i, vertexCount}] // Append[p]; tMax = Pi; trajectories = NDSolveValue[ Flatten @ { vertexRotationEquations, pivotFixtureEquation, vertexInitialConditions, pivotInitialCondition, collisionHandling }, dependentVariables, {t, 0, tMax} ]; vertexRealTrajectories[t_] := Most[trajectories][t] // Through // ReIm // Evaluate; Manipulate[ Show[ Graphics @ Polygon @ vertexRealTrajectories[tCurrent], If[tCurrent == 0, {}, ParametricPlot[ vertexRealTrajectories[t] , {t, 0, tCurrent} ] ] , ImageSize -> 240 ] , {tCurrent, 0, tMax} ] The nice thing about this method is that it generalises to any polygon. Also the initial pivot need not start on the bounding circle (and in fact need not be any of the vertices). At first we need to rotate the first edge around the first point.We use NMinimize to find such minimal angle. We ordering the points of polygon by clockwise,and rotate it by anticlockwise,so here we use RotateLeft to change the order of points. pts = {{1, 0}, {-0.2, -0.1}, {-0.1, 0.6}, {0.2, 0.8}, {Cos[1], Sin[1]}}; rotaAngle[p_, center_] := Module[{θ}, First@NMinimize[{θ, Norm@RotationTransform[θ, center][p] == 1, θ > 0}, {θ}]]; f[pts_] := Module[{θ, center}, θ = rotaAngle[pts[[2]], pts[[1]]]; center = pts[[1]]; RotateLeft[RotationTransform[θ, center] /@ pts, 1]]; polygon = NestList[f, pts, 20]; Manipulate[ Graphics[{Circle[], FaceForm[], EdgeForm[Red], Polygon[polygon[[n]]]}], {n, 1, 20, 1}] Here just a temporary workaround,I need time to find the way to rotate the polygon Independently. Edit Test another convex region. reg = ImplicitRegion[x^2/4 + y^2/3 <= 1 , {x, y}]; pts = {{2, 0}, {1, -0.2}, {1, 0.3}, {1.3, 0.6}, {1.5, 0.5}}; rotaAngle[pts_] := Module[{θ}, First@NMaximize[{θ, AllTrue[RotationTransform[θ, First[pts]] /@ Rest[pts], Element[#, reg] &], π > θ > 0}, {θ}, Method -> Automatic]]; f[pts_] := Module[{θ, center}, θ = rotaAngle[pts]; center = pts[[1]]; RotateLeft[RotationTransform[θ, center] /@ pts, 1]]; polygon = NestList[f, pts, 20]; regionplot = RegionPlot[reg, Frame -> False, PlotStyle -> {Cyan, Opacity[0.2]}, BoundaryStyle -> Green, AspectRatio -> Automatic]; Manipulate[ Show[regionplot, Graphics[{FaceForm[], EdgeForm[Red], Polygon[polygon[[n]]]}]], {n, 1, 20, 1}] • rotaAngle[p_, center_] := Module[{θ}, First@NMaximize[{θ, RotationTransform[θ, center][p] ∈ reg, π > θ > 0}, {θ}]]; Jan 9, 2021 at 0:14 A start at a better way: Graphics[ GeometricTransformation[Rectangle[], RotationTransform[30 Degree, {0, 0}]]]
2022-08-10T01:48:30
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https://www.jiskha.com/display.cgi?id=1303335098
# Calculus posted by . Gas is escaping a round balloon at2units^2/min.How fast is the surface area shrinking when the radius is 12units? • Calculus - Volume = V = (4/3) pi r^3 area =A = 4 pi r^2 r = (A/4pi)^.5 so V = (4/3) pi (A/4pi)^1.5 dV/dt = (4/3)(1/5)(A/4pi)^.5 (dA/dt/4pi) here dV/dt = -2 r = 12 so A = 4 pi(144) solve for dA/dt ## Similar Questions Related Rates: Gas is escaping from a spherical balloon at the rate of 2ft^3/min. How fast is the surface area shrinking (ds/dt) when the radius is 12ft? I posted this question a few days ago and someone answered it but I just wanted to make sure that I copied down the correct answers and didn't misunderstand anything. Related Rates: Gas is escaping a spherical balloon at the rate of … 3. ### calculus A balloon is being filled at a rate of 4000 cm^3/min. How fast is the surface area of the ballon increasing when the radius is 125 cm? 4. ### Calculus Air is being pumped into a spherical balloon so that its volume increases at a rate of 80 cm^3 per second. How fast is the surface area of the balloon increasing when its radius is 9cm? 5. ### math Example: 1 Gas is escaping from a spherical balloon at the rate of 2 ft3 / min .How fast is the surface area shrinking when the radius is 12 ft ? 6. ### calculus a spherical balloon filled with gas has a leak that permits the gas to escape at a rate of 1.5 cubic meters per minute. how fast is the surface area of the balloon shrinking when the radius is 4 meters? 7. ### Calculus: need clarification to where the #'s go Air is being pumped into a spherical balloon so that its volume increases at a rate of 80 \mbox{cm}^3\mbox{/s}. How fast is the surface area of the balloon increasing when its radius is 14 \mbox{cm}? 8. ### calculas i didnot under stand the question ? need help Gas is escaping a spherical balloon at the rate of 4 cm3 per minute. How fast is the surface area of the balloon shrinking when the radius of the balloon is 24 cm? 9. ### Math A spherical balloon expands uniformly as it is inflated. The radius of the balloon in r meters at time t seconds. Find the radius, the volume, and the surface area of the balloon when the rates of increase of the surface area of the … 10. ### Calculus A spherical balloon is being inflated and the radius of the balloon is increasing at a rate of 2cm/min. At what rates are the volume and surface area of the balloon increasing when the radius is 5cm? More Similar Questions
2017-08-23T02:37:34
{ "domain": "jiskha.com", "url": "https://www.jiskha.com/display.cgi?id=1303335098", "openwebmath_score": 0.9316695928573608, "openwebmath_perplexity": 476.40711070448714, "lm_name": "Qwen/Qwen-72B", "lm_label": "1. YES\n2. YES", "lm_q1_score": 0.9904405986777259, "lm_q2_score": 0.8459424353665381, "lm_q1q2_score": 0.8378557321313275 }
https://math.stackexchange.com/questions/3033045/how-to-find-the-matrix-represented-by-the-polynomials-a12-5a11-3i
# How to find the matrix represented by the polynomials $A^{12}-5A^{11}+…+3I$? I need to find the characteristic equation of the matrix $$A = \begin{bmatrix} 2&1&1\\ 0&1&0\\ 1&1&2\\ \end{bmatrix}$$ and find the matrix represented by $$A^{12}-5A^{11}+7A^{10}-3A^{9}-A^8+5A^7-7A^6+3A^5+A^4-5A^3+8A^2-4A+3I$$ My attempt: I started by finding the character equation which was $$\lambda^3-5\lambda^2+7\lambda-3=0$$ where $$\lambda$$ is the eigenvalue, and found $$\lambda = 5,\frac{9+\sqrt{69}}{2},\frac{9-\sqrt{69}}{2}$$ $$\lambda = 5$$ satisfies the equation and I would take $$\lambda = 5$$ let $$\lambda = A$$ and $$A^3-5 A^2+7A -3=0$$. I don't know how to proceed further. All help would be appreciated. Edit: What I can do, is multiply the matrix $$A$$ by itself and consequently the resultant matrix would be $$A^2$$ and so on. However, I'm not sure if this will give an accurate answer. • I suspect the characteristic equation is $\lambda^3 - 5\lambda^2 + 7\lambda - 3 = 0$. You might want to double check your work on that part. – JimmyK4542 Dec 9 '18 at 21:43 • Yes, you're right about the characteristic equation. I took a different problem's matrix for this one. My bad. – tNotr Dec 9 '18 at 21:48 The characteristic polynomial of $$A$$ is $$p(\lambda) = \det(\lambda I - A) = \lambda^3-5\lambda^2+7\lambda-3$$. So by the Cayley-Hamilton theorem $$p(A) = A^3-5A^2+7A-3I = 0$$ (*). Multiplying (*) by $$A^9$$ yields $$A^{12}-5A^{11}+7A^{10}-3A^9 = 0$$ (1). Multiplying (*) by $$-A^5$$ yields $$-A^8+5A^7-7A^6+3A^5 = 0$$ (2). Multiplying (*) by $$A$$ yields $$A^4-5A^3+7A^2-3A = 0$$ (3). If we add (1), (2), and (3), we get $$A^{12}-5A^{11}+7A^{10}-3A^9-A^8+5A^7-7A^6+3A^5+A^4-5A^3+7A^2-3A = 0$$ Can you take it from here? • I find ${}+8\lambda$ for the characteristic polynomial. – Bernard Dec 9 '18 at 21:53 • yes absolutely! However, I am not allowed to multiply the matrix A by any matrix greater than $A^3$ in the exam for some reason. Is this an attempt to baffle the student into making a mistake? I will never know but nonetheless I will try your solution though regardless of the risks. – tNotr Dec 9 '18 at 21:55 • You don't have to actually compute $A^3$ or any other higher power of $A$. Just add $A^2-A+3I$ to both sides of the last equation I wrote. Then, the left side is the expression you are looking for, and the right side is something that should be easy enough to compute. – JimmyK4542 Dec 9 '18 at 22:25 Check you computation for the characteristic polynomial $$\chi_A$$. Hint: Divide the polynomial $$p(x)=x^{12}-5x^{11}+\dots+3$$ by the characteristic polynomial: $$p(x)=q(x)\chi_A(x)+r(x)\qquad (\deg r \le 2),$$ to get $$p(A)=q(A)\chi_A(A)+r(A)=r(A).$$:
2019-01-16T09:45:59
{ "domain": "stackexchange.com", "url": "https://math.stackexchange.com/questions/3033045/how-to-find-the-matrix-represented-by-the-polynomials-a12-5a11-3i", "openwebmath_score": 0.8451789617538452, "openwebmath_perplexity": 194.55486529070296, "lm_name": "Qwen/Qwen-72B", "lm_label": "1. YES\n2. YES", "lm_q1_score": 0.958537730841905, "lm_q2_score": 0.8740772433654401, "lm_q1q2_score": 0.8378360174360565 }
https://math.stackexchange.com/questions/1664659/subring-proof-of-real-numbers
# subring proof of real numbers Problem: Show that the set of all real numbers of the form $$a_0 +a_1\pi +a_2\pi^2 + \cdot \cdot \cdot + a_n\pi^n$$ with $$n \ge 0, a_i \in \mathbb{Z}$$ is a subring of $$\mathbb{R}$$ that contains both $$\mathbb{Z}$$ and $$\pi$$. I think I've successfully shown all criteria of it being a subring. However, it's the statement about showing the it contains both $$\mathbb{Z}$$ and $$\pi$$ I'm not sure about. To show it's a subring must show it's closed under addition, multiplication, $$0_R$$ in the set, and the solution of $$a+x=1_R$$ is in the set. Let $$a_i,b_i\in \mathbb{Z}$$ with $$n\ge 0$$ then $$(a_0+a_1\pi +a_2\pi^2+\cdot \cdot \cdot + a_n\pi^n)+(b_0+b_1\pi+b_2\pi^2 + \cdot \cdot \cdot + b_n\pi^m$$ $$= (a_0+b_0)+(a_1+b_1)\pi+(a_2+b_2)\pi^2 + \cdot \cdot \cdot + (a_n+b_n)\pi^{n+m}$$ Similarly, $$(a_0+a_1\pi +a_2\pi^2+\cdot \cdot \cdot + a_n\pi^n)(b_0+b_1\pi+b_2\pi^2 + \cdot \cdot \cdot + b_n\pi^m)$$ $$= (a_0b_0)+(a_0b_1+a_1b_0)\pi+(a_0b_2+a_1b_1+a_2b_0)\pi^2 + \cdot \cdot \cdot + (a_nb_n)\pi^{n+m}$$ Both of which have the required form by polynomial addition and multiplication. $$0_R = a_0 +a_1\pi+a_2\pi^2 + \cdot \cdot \cdot +a_n\pi^n$$ where $$(\forall i\in (0,n))(a_i = 0)$$ $$-(a_0 +a_1\pi + \cdot \cdot \cdot + a_n\pi^n)$$ is in the set since $$-a_i\in \mathbb{Z}$$, hence the solution to $$a+x=0_R$$ is in the set. Now, for the last part (assuming all of the above is right showing that its a subring). Is it enough to just state that since each $$a_i \in \mathbb{Z}$$ it must necessarily contain $$\mathbb{Z}$$ since each $$a_i$$ is arbitrary. Each $$a_i$$ with a coefficient of $$\pi$$ isn't in $$\mathbb{Z}$$ since an integer times an irrational is still irrational, but since $$a_0$$ is just a constant term with an arbitrary $$a_i\in \mathbb{Z}$$ it can span the integers and hence the subring contains $$\mathbb{Z}$$. For the last part it obviously contains multiples of $$\pi$$ • It is enough to state that for each $m\in\mathbb{Z}$, setting $a_1=a_2=\ldots=0$ and $a_0=m$ retrieves the element $m$. The notation "$R$ contains $\mathbb{Z}$ and $\pi$" is a bit dangerous. I assume they mean $\{\pi\}$. Feb 20 '16 at 20:57 Now, $\mathbf{Z}$ is contained in out set because the first coefficient $a_0$ is an integer, and we can set all other coefficients to zero. Certainly $\pi$ is in our set because we can set all coefficients other than $a_1$ to zero, and let $a_1 = 1$. So you're done, your proof seems flawless to me.
2021-12-03T09:43:58
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https://math.stackexchange.com/questions/209422/is-the-variance-concave
# Is the variance concave? Let $$X$$ be a discrete random variables with values in the set $$\{x_1,\ldots, x_n\}\subset\mathbb{R}$$. Denote by $$p_i$$ the probability that $$X=x_i$$. We can then regard the variance $$\mbox{Var}(X)$$ as a function of the vector $$p\in \Delta^{n-1}\subseteq\mathbb{R}^n$$. Will it be a concave function of $$p$$? With $$n=2$$, we get $$\mbox{Var} (X)=p_1p_2(x_1-x_2)^2=p_1(1-p_1)(x_1-x_2)^2$$ which is concave. I'm not sure how to generalize this beyond two dimensions though. We have $\mathrm{var}_p\left(X\right)=\mathbb{E}_p\left[X^2\right]-\mathbb{E}_p\left[X\right]^2$. Furthermore, for $\lambda\in\left[0,1\right]$ we can write $$\mathbb{E}_{\lambda p+(1-\lambda)q}\left[X^2\right]=\lambda\mathbb{E}_p\left[X^2\right]+(1-\lambda)\mathbb{E}_q\left[X^2\right].$$ Convexity of $\left(\cdot\right)^2$ and Jensen's inequality then yield $$\lambda \mathbb{E}_p\left[X\right]^2+(1-\lambda) \mathbb{E}_q\left[X\right]^2\geq \left(\lambda\mathbb{E}_p\left[X\right]+(1-\lambda)\mathbb{E}_q\left[X\right]\right)^2=\mathbb{E}_{\lambda p+(1-\lambda)q}\left[X\right]^2.$$ Thus we deduce $$\mathrm{var}_{\lambda p+(1-\lambda)q}\left(X\right)=\mathbb{E}_{\lambda p+(1-\lambda)q}\left[X^2\right]-\mathbb{E}_{\lambda p+(1-\lambda)q}\left[X\right]^2\\\geq\lambda\mathrm{var}_p\left(X\right)+(1-\lambda)\mathrm{var}_q\left(X\right),$$ that means the variance is concave. • Thank you, that was a very clean solution. Oct 8, 2012 at 20:03 We look at the Hessian and see it's negative semi-definite because i) function itself \begin{align*} f(p) &= \text{Var}(X)\\ &= E[X^2] - E[X]^2\\ &= p^Ty - (p^Tx)^2, \end{align*} where $$y = [x_1^2,\,\ldots,\,x_n^2]^T$$. \begin{align*} \nabla f &= y - 2(p^Tx)x\\ \nabla f^2 &= -2xx^T. \end{align*} $$\nabla f^2$$ is negative semi-definite for all $$p$$ $$\Leftrightarrow$$ $$f$$ is concave.
2022-05-24T16:03:23
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https://math.stackexchange.com/questions/421099/fracc1c-leq-fraca1a-fracb1b-for-0-leq-c-leq-ab
# $\frac{c}{1+c}\leq\frac{a}{1+a}+\frac{b}{1+b}$ for $0\leq c\leq a+b$ When proving that if $d$ is a metric, then $d'(x,y)=\dfrac{d(x,y)}{1+d(x,y)}$ is also a metric, I have to prove the inequality: $$\dfrac{c}{1+c}\leq\frac{a}{1+a}+\frac{b}{1+b}$$ for $0\leq c\leq a+b$. This is obvious by expanding, but is there a nicer way to see why it is true? • You can also use the monotone increasing property of $\frac{x}{1+x}$ for proving $\frac{d(x,y)}{1+d(x,y)}$ is a metric. – mtm Jun 17 '13 at 17:50 $$\frac{a}{1+a}+\frac{b}{1+b}\ge\frac{a}{1+a+b}+\frac{b}{1+b+a}=\frac{a+b}{1+a+b}=\frac{1}{1+\frac{1}{a+b}}\ge\frac{1}{1+\frac1c}=\frac{c}{1+c}$$ • this comment was dumb, sorry, really need to get some sleep. cheers – mm-aops Jun 15 '13 at 13:23 • Very clear, thanks Maisam! – Paul S. Jun 15 '13 at 13:27 • @mm-aops:no problem care free – M.H Jun 15 '13 at 13:28 • @Paul S.:your welcome dear paul – M.H Jun 15 '13 at 13:29 $\displaystyle\frac{d(x,z)}{1+d(x,z)}-\frac{d(x,y)}{1+d(x,y)}=\frac{(d(x,z)-d(x,y))}{(1+(d(x,z))(1+d(x,y))}$ We havefrom $\Delta$ inequality, $(1+(d(x,z))(1+d(x,y))\ge1+d(x,z)+d(x,y)\ge 1+d(y,z)\tag 1$ And we have from $\Delta$ inequality $d(x,z)-d(x,y)\le d(y,z)\tag 2$ And as $1+d(y,z)\ne 0$ we have $\displaystyle\frac{1}{1+d(y,z)}\ge \frac{1}{(1+d(x,z))(1+d(x,y))}$ from $(1)$ Multiplying $(2)$ and $(1)$ we have $$\frac{d(x,z)}{1+d(x,z)}-\frac{d(x,y)}{1+d(x,y)}=\frac{(d(x,z)-d(x,y))}{(1+(d(x,z))(1+d(x,y))}\le \frac{d(y,z)}{1+d(y,z)}$$ $$\Rightarrow \frac{d(x,z)}{1+d(x,z)}-\frac{d(x,y)}{1+d(x,y)}\le \frac{d(y,z)}{1+d(y,z)}$$ $$\Rightarrow \frac{d(x,z)}{1+d(x,z)}\le \frac{d(x,y)}{1+d(x,y)}+ \frac{d(y,z)}{1+d(y,z)}$$
2019-09-21T17:30:35
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https://math.stackexchange.com/questions/2946288/higher-dimensional-real-matrix-representation-of-imaginary-unit/2946485
# Higher dimensional real matrix representation of imaginary unit? The imaginary unit can be represented as real 2x2 matrix : $$i = \begin{pmatrix} 0 & -1 \\ 1 & 0\\ \end{pmatrix},\, i^2 = \begin{pmatrix} -1 & 0 \\ 0 & -1\\ \end{pmatrix}$$ (see https://en.wikipedia.org/wiki/Imaginary_unit#Matrices). Now from my shallow knowledge of group representation theory I believe that it should also be possible to find higher dimensional real representations of $$i$$. I tried to find a 3x3 representation, which is basically asking if there is matrix square root like $$\sqrt{\begin{pmatrix} -1 & 0 & 0 \\ 0 & -1 & 0\\0 & 0 & -1 \end{pmatrix}} \sqrt{\begin{pmatrix} -1 & 0 & 0 \\ 0 & -1 & 0\\0 & 0 & -1 \end{pmatrix}} = \begin{pmatrix} -1 & 0 & 0 \\ 0 & -1 & 0\\0 & 0 & -1 \end{pmatrix}$$ but couldn't find any (by try and error). Nevertheless I was able to find a 4x4 matrix which can be used as representation of $$i$$: $$\begin{pmatrix} 0 & -1 & 0 & 0\\ 1 & 0 & 0 & 0\\0& 0 & 0 & -1\\ 0 & 0 & 1 & 0 \end{pmatrix}\begin{pmatrix} 0 & -1 & 0 & 0\\ 1 & 0 & 0 & 0\\0& 0 & 0 & -1\\ 0 & 0 & 1 & 0 \end{pmatrix} = \begin{pmatrix} -1 & 0 & 0 & 0 \\ 0 & -1 & 0 & 0\\0 & 0 & -1 & 0 \\ 0 & 0 & 0 & -1 \end{pmatrix}$$ So can it be that $$i$$ can only be represented as real matrix of even dimensions $$N=2n\, (n>0)$$? Can this be proven and maybe a construction be given on how to construct higher dimensional real matrix representations of the imaginary unit? And can it be shown, that the 2x2 representation is really the lowest dimensional representation of $$i$$ as a real matrix? • If $A^2=-I_3$, then $(\det{A})^2 = -1$. – Catalin Zara Oct 7 '18 at 20:45 • @CatalinZara Why not an official answer? – Paul Frost Oct 7 '18 at 22:56 • @CatalinZara Ok, from math.stackexchange.com/questions/2408343/… I finally understand what you mean: "In odd dimensions however, there is no real matrices which are skew-symmetric and orthogonal. " – asmaier Oct 8 '18 at 20:35 Suppose $$T$$ is an $$n\times n$$ matrix such that $$T^2=-I$$. Note then then $$\mathbb{R}^n$$ becomes a vector space over $$\mathbb{C}$$ (extending the $$\mathbb{R}$$-vector space structure), by letting $$a+ib\in\mathbb{C}$$ act by $$aI+bT$$. But then by picking a basis for this vector space over $$\mathbb{C}$$, we have $$\mathbb{R}^n\cong \mathbb{C}^m$$ for some $$m$$. As a real vector space, $$\mathbb{C}^m\cong \mathbb{R}^{2m}$$, and so it follows that $$n=2m$$ is even. Conversely, for any $$m$$, $$\mathbb{C}^m$$ is a $$2m$$-dimensional real vector space, and multiplication by $$i$$ is an $$\mathbb{R}$$-linear map whose square is $$-1$$. So such a matrix does exist whenever $$n$$ is even. Explicitly, this is just the obvious generalization of your $$4\times 4$$ matrix, a block diagonal matrix with the $$2\times 2$$ diagonal matrix for $$i$$ on the diagonal blocks. I hope Catalin Zara's comment has convinced you that you can only find a $$n\times n$$ matrix $$M$$ with $$M^2=-1$$, if $$n=2k$$. Here I will construct such $$M$$ in full generality. Let us make a few definitions. $$\Sigma_{(p,q)} = \begin{pmatrix} 1_{p\times p} &0\\ 0 & -1_{q\times q} \end{pmatrix}, \qquad Y=\begin{pmatrix} 0 & -1\\ 1 & 0 \end{pmatrix}$$ where $$\Sigma_{(p,q)}$$ is $$(p+q)\times (p+q)$$ , while $$Y$$ is $$2\times 2$$. Now, if you use Jordan normal form for real matrices, then with a little bit of work, you find that if $$M$$ satisfies $$M^2=-1$$, then its real Jordan normal form is necessarily of the form $$J = \Sigma_{(p,q)}\otimes Y$$ where $$p+q=k$$ and $$\otimes$$ is the Kronecker product. As a result, $$M^2=-1$$ if and only if $$M = P^{-1}(\Sigma_{(p,k-p)}\otimes Y) P$$ where $$P$$ is an invertible real matrix and $$0\leq p\leq k$$. You can think of $$p$$ as "signature" of $$M$$, which classifies matrices with $$M^2=-1$$, up to a similarity relation. • Could you elaborate on Catalin Zaras comment, because I didn‘t understand that. – asmaier Oct 8 '18 at 8:04
2019-09-18T22:00:15
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https://math.stackexchange.com/questions/3826019/how-many-values-of-x-in-mathbb-z-x99-are-there-such-that-m-n-in-mathbb-z/3826040
# How many values of $x\in\mathbb Z^+,x<99$ are there such that $m,n\in\mathbb Z$ and $m^2-n^2=x$ is possible? How many values of $$x\in\mathbb Z^+,x<99$$ are there such that $$m,n\in\mathbb Z$$ and $$m^2-n^2=x$$ is possible? So what I'm trying to find here is the number of integers between $$1$$ and $$98$$ inclusive such that that integer can be expressed as the difference of two squares. I know that all odd numbers can be expressed as the difference between to consecutive squares, so the answer is at least $$98/2=49$$, but I don't really see a way to continue from here. Maybe I can utilize Pythagorean Theorem somehow? Thanks for the help. Also I'm not too sure which topic this question falls under so if someone could edit the tags that would be great. • $m^2-n^2=(m-n)(m+n)$ – player3236 Sep 14 '20 at 16:48 • @player3236 Yes I have noticed that but I'm not sure how I can utilize this information to solve the problem. – Aiden Chow Sep 14 '20 at 16:49 • Just take the differences of squares of moderate size. So $3=2^2-1^2$ is possible, etc. You can use $x=(m+n)(m-n)$ to derive a contradiction for certain $x$, like $x=2$. – Dietrich Burde Sep 14 '20 at 16:50 • @DietrichBurde I'm still not seeing it. How do I know which values of $x$ to test, and how exactly do you get a contradiction from $(m+n)(m-n)$ for certain values of $x$? – Aiden Chow Sep 14 '20 at 16:53 • Indeed these are the only ones. – player3236 Sep 14 '20 at 17:00 There are two ways to generate an even number as the difference of two squares: $$m,n$$ are either both even or both odd. If $$m=2k$$ and $$n=2l$$, then $$m^2-n^2=4(k^2-l^2)$$. This immediately tells us that all numbers $$4x$$, $$x$$ odd, are possible. If $$m=2k+1$$ and $$n=2l+1$$: $$m^2-n^2=4(k^2-l^2+k-l)=4(k-l)(k+l+1)$$ and every even number $$2x$$ may be written in the form $$(k-l)(k+l+1)$$ by setting $$k=x$$ and $$l=x-1$$. So all numbers $$4x$$, $$x$$ even, are possible. In conclusion, the numbers that are the difference of two squares are odd numbers and multiples of $$4$$. There are $$49+\lfloor98/4\rfloor=49+24=73$$ such numbers in the given range. I'll supplement the other proof by showing also that numbers that are even but not divisible by $$4$$ cannot be expressed as a difference of squares of integers. Suppose $$x = m^2-n^2 = (m-n)(m+n)$$. Let $$p,q\in\mathbb Z$$ such that $$x = pq$$. By letting $$p = m-n$$, $$q=m+n$$, we can express $$x$$ as a difference of squares. However this requires $$p+q = 2m, q-p = 2n$$ to both be even. This is acheived only if both $$p,q$$ are even or odd. If both $$p,q$$ are even, $$x$$ is divisible by $$4$$. If both $$p,q$$ are odd, $$x$$ is also odd. If $$x$$ is even but not divisible by $$4$$, it must be a product of an odd number and an even number. Then $$p+q$$ and $$q-p$$ must be odd, which does not fulfil our requirement. This method is also constructive: any factorization of $$x$$ gives a solution to the difference of squares. watch this! If $$x=2h+1$$ is odd then $$x = 1*(2h+1)= ([h+1]-h)([h+1]+h) =(h+1)^2 - h^2$$ so every odd number is difference of square. If we want to be a bit more creative if $$x = j*k$$ an odd composite then $$j=\frac {j+k}2 + {j-k}2$$ and $$k = \frac {j+k}2-\frac {j-k}2$$ so $$x = (m+n)(m-n) = m^2 - n^2$$ if $$m=\frac {j+k}2$$ (which is an integer as both $$j,k$$ are odd) and $$n = \frac {j-k}2$$ (ditot). So every odd number will work. If $$x = 2w$$ and $$w$$ is odd, then if $$x = m^2 -n^2 = (m-n)(m+n)$$ then one of $$m+n$$ or $$n-m$$ is even and the other is not. $$m+n = (m-n) + 2n$$ so if $$m-n$$ is even or odd then so is $$m+n =(m-n)+2n$$. so that is impossible. So every number that is even but not divisible by $$4$$ will not work. Induction time! Now if $$w= m^2- n^2$$ is possible, and $$x= 4w$$ then $$4x = (2m)^2 - (2n)^2$$ is will work. And note if $$w = m^2 - n^2= (m-n)(m+n)$$ then $$8w = 2*4(m-n)(m+n)=(2m-2n)(4m+4m) = [(3m+n)-(m+3n)][(3m+n) + (m+3n)] = (3m+n)^2 - (m+3n)^2$$ will work. So by induction, if $$w=m^2 - n^2$$ will work then $$2^kw$$ will work if $$k$$ is even, or if $$k$$ is a multiple of $$3$$ or $$k$$ is a sum of an even number and a multiple of $$3$$. But that can be any positive integer greater than $$1$$. And as $$h$$ odd will work, and $$2h; h$$ odd will not, then $$2^kh=4(2^{k-1} h); h$$ odd$$; k\ge 2$$ will work. ===== So the numbers that can be so written are: Every odd number, every multiple of $$4$$, but no number that is even but not divisible by $$4$$. Using Eudlid's formula $$\quad A=m^2-k^2,\quad B=2mk,\quad C=m^2+k^2\quad$$ or any other, we know that the value of $$A$$ can be any odd number except $$1$$, i.e. $$A\ge3$$. This means that, under $$99$$, there are $$98/2-1=48$$ odd values of $$A$$. It also happens that, with this formula ––the only one that uses $$m^2-k^2$$–– $$A$$ can be any multiple of $$4$$ greater than $$4$$ such as $$f(3,1)=(8,6,10)\quad f(4,2)=(12,16,20)\quad f(5,3)=(16,30,34)\quad...\quad f(10,2)=(96,40,104)$$ Note $$\quad f(2,0)=(4,0,4)\qquad$$ but this is a trivial triple so most people do not treat it as valid. This means that there are $$96/4-1=24-1=23$$ even values of $$A$$ that can be generated by $$m^2-k^2$$ for $$A<99$$. The total is $$48+23=71$$.
2021-01-19T05:21:22
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https://math.stackexchange.com/questions/1240279/how-to-find-the-maximum-and-minimum-of-the-function-fx-frac3xx2-2x/1240288
# How to find the maximum and minimum of the function $f(x) = \frac{3x}{x^2 -2x + 4}$ How would one find the maximum and minimum of such a function: $$f: \mathbb{R} \rightarrow \mathbb{R}, x \mapsto f(x) = \frac{3x}{x^2 -2x + 4}$$ I have just been introduced to functions in my calculus class (actually I missed that lectures), and I have more or less have the intuition of maximum and minimum of a function: the maximum should be the greatest $f(x)$ value in a certain range and the minimum the smallest $f(x)$ value. My first question is: does it makes sense to talk about maximum and minimum of a function in general (instead of just in a certain subset of the domain of the function)? I need to use maxima to find the maximum and minimum, and then I need to show formally that they are really the maximum and minimum. My second question is: does anybody know how to find the maximum and minimum of a function using wxMaxima? My third question is: I know the definition of maximum (and minimum), i.e. a number that is greater or equal (smaller or equal) to all other numbers. My problem is that I am not seeing how would I show it formally. I used Wolfram Alpha to calculate the maximum and the minimum, and it says that the minimum is $-\frac{1}{2}$ at $-2$ and the maximum is $\frac{3}{2}$ at $2$. If one lets $y=\frac{3x}{x^2-2x+4},$ one has $y=0\iff x=0$. Suppose that $y\not=0$. Since $y(x^2-2x+4)=3x\iff yx^2+(-2y-3)x+4y=0$, considering the discriminant gives you $$(-2y-3)^2-4\cdot y\cdot (4y)\ge 0\iff (2y-3)(2y+1)\le 0.$$ Hence, one has $-\frac 12\le f(x)\le \frac 32$. Here, the equalities are attained : $f(-2)=-\frac 12,f(2)=\frac 32$. • Why should this give us the maximum and minimum? I am not seeing the connection... – nbro Apr 18 '15 at 12:33 • @Rinzler: Sorry, but please clarify what you don't understand in my answer. Do you understand $-\frac 12\le f(x)\le \frac 32$? – mathlove Apr 18 '15 at 12:38 • @Rinzler: (1) The reason why I separate it into two cases $y=0$ and $y\not=0$ is that I wanted to use the discriminant for $x$. Note that the coefficient of $x^2$ is $y$. This must not be $0$. (2) Since $x$ is real, the discriminant has to be non-negative (I'm sure you know this) (3) When $y\not=0$, I got $(2y-3)(2y+1)\le 0$. Combining this with $y=0$, I got $-\frac 12\le y\le \frac 32$. (4) Considering the discriminant sometimes works to find the maximum and the minimum because it gives an inequality for $y$. – mathlove Apr 18 '15 at 12:48 • nice work. i would have done the same way. @Rinzler, read the answer more carefully. – abel Apr 18 '15 at 13:25 • @Rinzler: You can get it by expanding LHS : $(-2y-3)^2-4\cdot y\cdot (4y)\ge 0\iff 4y^2+12y+9-16y^2\ge 0\iff -12y^2+12y+9\ge 0$ $\iff -4y^2+4y+3\ge 0\iff 4y^2-4y-3\le 0\iff (2y-3)(2y+1)\le 0$. – mathlove Apr 19 '15 at 22:04 solve the equation $f'(x)=-\frac{3 (x-2) (x+2)}{\left(x^2-2 x+4\right)^2}=0$ for $x$. if you know it is $\frac{3}{2}$ you can calculate $$\frac{3}{2}-\frac{3x}{x^2-2x+4}=\frac{3(x-2)^2}{2(4-2x+x^2)}\geq 0$$ • We have not been introduced to the concept of derivatives... – nbro Apr 18 '15 at 12:11 Since the discriminant of $x^2-2x+4$ is negative we have that $x^2-2x+4$ never vanishes, so it is enough to compute the minimum of $g(x)=\frac{1}{f(x)}$ over $\mathbb{R}^+$ and the maximum of $\frac{1}{f(x)}$ over $\mathbb{R}^-$, since $x=0$ is the only zero of $f(x)$ and $\lim_{x\to\pm\infty}f(x)=0$. Now: $$g(x) = \frac{x}{3}+\frac{4}{3x}-\frac{2}{3},$$ so the stationary points of $g(x)$, by the AM-GM inequality, occur when $\frac{x}{3}=\frac{4}{3x}$, i.e. at $x=\pm 2$. That directly gives: $$-\frac{1}{2}= f(-2) \leq f(x) \leq f(2) = \frac{3}{2}$$ as wanted. • What do you mean by: "$x^2-2x+4$ never vanishes"? Why is it enough to computer the minimum of $g(x)$ over $R^+$ and the maximum over $R^{-1}$? What do you mean by "over $R^{+/-}$"?..... – nbro Apr 19 '15 at 14:02 • @Rinzler: "never vanishes" means that for every $x\in\mathbb{R}$ we have $x^2-2x+4\neq 0$. $\mathbb{R}$ is the set of real numbers, $\mathbb{R}^+$ the set of positive real numbers, $\mathbb{R}^-$ the set of negative real numbers. If a continuous function $g$ never vanishes, the maxima of $g$ are attained in the minima of $\frac{1}{g}$ and vice-versa. – Jack D'Aurizio Apr 19 '15 at 14:19 If $x\ne0,$ $$\frac{3x}{x^2-2x+4}=\dfrac3{x-2+4/x}$$ If $x>0,x-2+4/x=\left(\sqrt x-2/\sqrt x\right)^2+4-4\ge2$ Consequently, $$\frac{3x}{x^2-2x+4}\le\dfrac32$$ If $x<0,$ let $x=-y^2,x-2+4/x=-y^2-2-4/y^2=-2-\left(y-2/y\right)^2-4\ge-6$ Consequently, $$\frac{3x}{x^2-2x+4}\ge\dfrac3{-6}=-\dfrac12$$
2020-10-29T05:56:45
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https://math.stackexchange.com/questions/3732792/probabilities-in-the-expectation-of-a-hypergeometric-random-variable
# Probabilities in the expectation of a hypergeometric random variable So I've read this explanation and the answer given by Andi R helped to give me some intuition of my problem. Now would like to understand it more mathematically rigorously. Suppose we have a jar of marbles with $$r$$ red balls and $$w$$ white balls. We draw n balls. Find the expected number of red balls. From the expectation of a hypergeometric random variable know that the expected number of red balls is going to be $$\frac{nr}{r + w}$$ Let's now use the method of indicators random variables to point at my problem. Let $$X$$ be the number of red balls. $$X = X_1 + X_2 + X_3 +... + X_n$$ \begin{align*} X_i&=\begin{cases} 1,&\text{if }i\text{-th ball selected is red}\\ 0,&\text{otherwise} \end{cases}\\ \end{align*} By linearity of expectation, $$E[X] = \sum_i{E[X_i]} = \sum_i {P(X_i = 1)}$$ That all makes sense to me. Here's where I start to get confused. $$P(X_i = 1) = \frac{r}{r+w}\quad \forall i$$ However let's look at our second ball for example, using the law of total probability: $$P(X_2 = 1) = P(X_2 = 1| X_1 = 1)P(X_1 = 1) + P(X_2 = 1| X_1 = 0)P(X_1 = 0)$$ $$P(X_2 = 1) = \frac{r - 1}{r + w - 1}\frac{r}{r+w} + \frac{r}{r + w - 1}\frac{w}{r+w} \neq \frac{r}{r+w}$$ Where am I going wrong? Did I use the total law of probability wrong? We draw the balls successively so the probabilities can't be the same, can it? • I can only find one mistake which (fortunately) is $\cdots\neq\cdots$. Jun 24, 2020 at 13:58 • $(r-1)r+rw=r(r+w-1)$ Jun 24, 2020 at 13:59 • You two are totally correct, I made a algebraic error on my scrap paper! Jun 24, 2020 at 14:01 • Thank you so much! Jun 24, 2020 at 14:02 $$P(X_2 = 1) = \frac{r - 1}{r + w - 1}\frac{r}{r+w} + \frac{r}{r + w - 1}\frac{w}{r+w} \\ = \frac{r(r-1)}{(r+w-1)(r+w)} + \frac{rw}{(r+w-1)(r+w)} \\ = \frac{r(r-1) + rw}{(r+w-1)(r+w)} \\ = \frac{r^2 +rw-r}{(r+w-1)(r+w)}\\ = \frac{r(r + w-1)}{(r+w-1)(r+w)}\\ = \frac{r}{r+w}$$ And a very similar algebraic motions can be followed for $$X_3 = 1$$, $$X_4 = 1$$, and so forth...
2022-08-10T06:41:58
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https://math.stackexchange.com/questions/1952056/computing-the-integral-int-expix2-dx
# Computing the integral $\int \exp(ix^2) dx$ I'm trying to compute the following integral: $I_1 = \int_{-\infty}^\infty \exp(iu^2) du$. This is what I did, but it is wrong and I don't know why: $$I_1^2 = \left (\int_{-\infty}^\infty \exp(iu^2) du \right)^2 = \int_{-\infty}^{\infty} \exp(iu^2) du \int_{-\infty}^{\infty} \exp(iv^2) dv = \iint_{\mathbb{R}^2} \exp{i(v^2+y^2)} dA$$ Using polar coordinates: $$I_1^2 = \lim_{c\to \infty} \int_0^{2\pi} d\varphi \int_0^{c} \exp(ir^2) r dr$$ But the last integral doesn't converge. What is wrong here? • Is there a reason to believe that your original integral converges? After all, the integrand has constant magnitude $|e^{iu^2}|=1$. – MPW Oct 3 '16 at 13:34 • @MPW, It is not hard to prove that it converges conditionally. – Sangchul Lee Oct 3 '16 at 13:36 • Check out Fresnel Integrals here en.wikipedia.org/wiki/Fresnel_integral – Kevin Oct 3 '16 at 13:49 • Of course re-arranging a conditionally convergent integral does not work ... such as converting from cartesian to polar coordinates. – GEdgar Oct 3 '16 at 13:50 • @GEdgar why it does not work? That is what Sangchul Lee did. – Tom Oct 3 '16 at 14:28 In the last equality of the second equation, you tried to utilize Fubini's theorem. But the issue is that the integral is only conditionally convergent. So in principle, you cannot apply Fubini's theorem to convert your iterated improper integral into the double integral. Assuming that you have already established the existence of $I$, here are two workarounds that savages your approach. Notice that both methods approximate the original integral by absolutely convergent ones so that Fubini's theorem is applicable. Method 1. Consider the truncated version $$I_R = \int_{-R}^{R} e^{ix^2} \, dx.$$ Then by Fubini's theorem, we have $$I_R^2 = \iint_{[-R,R]^2} e^{i(x^2+y^2)} \, dxdy.$$ Divide the square region $[-R,R]^2$ into 4 congruent pieces along lines $y = x$ and $y = -x$. This will result in 4 integrals. Exploiting the symmetry, it is easy to check that these 4 integrals have identical values, so you can also write: $$I_R^2 = 4 \iint\limits_{\substack{0 \leq x \leq R \\ |y| \leq R}} e^{i(x^2+y^2)} \, dxdy.$$ Now converting this integral using polar coordinates, we get $$I_R^2 = 4 \int_{-\frac{\pi}{4}}^{\frac{\pi}{4}} \int_{0}^{R\sec\theta} e^{ir^2} \, rdrd\theta = 2i \int_{-\frac{\pi}{4}}^{\frac{\pi}{4}} (1 - e^{iR^2\sec^2\theta}) \, d\theta.$$ As $R \to\infty$, Riemann-Lebesgue lemma shows that this integral converges to $$I^2 = \lim_{R\to\infty} I_R^2 = 2i \int_{-\frac{\pi}{4}}^{\frac{\pi}{4}} \, d\theta = i\pi.$$ This determines $I$ up to sign, but it is not hard to check that $\operatorname{Im}(I) > 0$. Thus the only possible choice is $$I = \sqrt{i\pi} = (1+i)\sqrt{\frac{\pi}{2}}.$$ Method 2. Alternatively, you can use Gaussian regularization: $$I = \lim_{\epsilon \downarrow 0} I_{\epsilon} \quad \text{where} \quad I_{\epsilon} = \int_{-\infty}^{\infty} e^{ix^2}e^{-\epsilon x^2} \, dx.$$ Then $I_{\epsilon}^2$ is absolutely convergent and Fubini's therorem followed by polar coordinates change gives $$I_{\epsilon}^2 = \int_{0}^{2\pi} \int_{0}^{\infty} e^{-(\epsilon-i)r^2} \, rdrd\theta = \frac{\pi}{\epsilon - i}.$$ This converges to $\pi i$ as $\epsilon \downarrow 0$ and you can proceed from here as in Method 1. Existence of $I$. So here is one way to establish the existence of $I$. A moment of thought tells us that it suffices to show the convergence of $\int_{1}^{\infty} e^{ix^2} \, dx$. Let $R > 1$. The usual trick is to improve the speed of convergence: $$\int_{1}^{R} e^{ix^2} \, dx \ \stackrel{u=x^2}{=} \int_{1}^{R^2} \frac{e^{iu}}{2\sqrt{u}} \, du.$$ Applying integration by parts, $$\int_{1}^{R} e^{ix^2} \, dx = \left[ \frac{e^{iu}}{2i\sqrt{u}} \right]_{1}^{R} + \int_{1}^{R} \frac{e^{iu}}{4iu^{3/2}} \, du.$$ Now the integral on the right-hand side is absolutely convergent and we are done. Another method is as follows: Write $$\int_{1}^{R} e^{ix^2} \, dx = \bigg( \sum_{k = 0}^{\lfloor R^2/\pi \rfloor - 1} \int_{\sqrt{\pi k}}^{\sqrt{\smash[b]{\pi(k+1)}}} e^{ix^2} \, dx \bigg) + \int_{\sqrt{\pi \smash[b]{\lfloor R^2/\pi \rfloor}}}^{R} e^{ix^2} \, dx.$$ Then you can show that the sum converges by applying the alternating series test both to the real part and to the imaginary part. The error term can be shown to vanish as $R \to \infty$ as well. Here we prove the following claim: Proposition. Suppose that $f : \Bbb{R} \to \Bbb{C}$ is locally integrable on $\Bbb{R}$ and is improperly integrable on $\Bbb{R}$ in the sense that the limit $$I(0) := \lim_{\substack{a \to -\infty \\ b \to \infty}} \int_{a}^{b} f(x) \, dx$$ converges. Then for each $\epsilon > 0$, the integral $$I(\epsilon) := \int_{-\infty}^{\infty} f(x)e^{-\epsilon x^2} \, dx$$ is well-defined and $I(0) = \lim_{\epsilon \to 0^+} I(\epsilon)$. The proof is quite standard and easy. Let $F : \Bbb{R} \to \Bbb{C}$ be defined by $$F(x) = \int_{0}^{x} f(t) \, dt.$$ Then the assumptions tells us the following things: • Both $F(\infty) := \lim_{b\to\infty} F(b)$ and $F(-\infty) := \lim_{a\to-\infty} F(a)$ converge, • $F$ is bounded, and • $I(0) = F(\infty) - F(-\infty)$. Now, from integration by parts, $$\int_{a}^{b} f(x) e^{-\epsilon x^2} \, dx = \left[ F(x)e^{-\epsilon x^2} \right]_{a}^{b} + \int_{a}^{b} 2\epsilon x F(x) e^{-\epsilon x^2} \, dx.$$ Taking $(a, b) \to (-\infty, \infty)$, dominated convergence theorem shows the above integral converges and \begin{align*} I(\epsilon) = \int_{-\infty}^{\infty} f(x) e^{-\epsilon x^2} \, dx &= \int_{-\infty}^{\infty} 2\epsilon x F(x) e^{-\epsilon x^2} \, dx \\ &= \int_{-\infty}^{\infty} 2u F(\epsilon^{-1/2}u) e^{-u^2} \, du, \end{align*} where the substitution $u = \sqrt{\epsilon}x$ is used in the last line. Now we focus on the last integral and notice that the integrand is dominated by $C|u|e^{-u^2}$ for some absolute constant $C \geq 0$, which is integrable. (Here, $C = 2 \sup_{x\in\Bbb{R}} |F(x)| < \infty$ works for the proof.) Therefore, by the dominated convergence theorem, \begin{align*} \lim_{\epsilon \to 0^+} I(\epsilon) &= \int_{-\infty}^{\infty} \lim_{\epsilon \to 0^+} 2u F(\epsilon^{-1/2}u) e^{-u^2} \, du \\ &= \int_{-\infty}^{0} 2u F(-\infty) e^{-u^2} \, du + \int_{0}^{\infty} 2u F(\infty) e^{-u^2} \, du \\ &= F(\infty) - F(-\infty) = I(0). \end{align*} • Great answer :) I got it all, except the Gaussian regularization. Where can I read about it? This is the first time I hear about that concept. And how can I prove that is only conditionally convergent? And how can I establishe the existence of $I$? – Tom Oct 3 '16 at 14:15 • @Tom, It is just a specific way of approximating an integral by nicer ones. I guess you will find similar ones in many real analysis books, and I can also prove the specific version above for you if you want. – Sangchul Lee Oct 3 '16 at 14:20 • That would be great :) – Tom Oct 3 '16 at 14:22 • @Tom, I have to leave now and will explain that this afternoon, so please remind me if I forget to do so within today :) – Sangchul Lee Oct 3 '16 at 14:43 • I got it, thanks :) Where can I read about Gaussian regularization? your 2nd method. – Tom Oct 3 '16 at 22:17 This and related ones are known as Fresnel Integrals, see HERE to find that (improper Riemann integrals) $$\int_{-\infty}^\infty e^{ix^2}dx = \int_{-\infty}^\infty \cos(x^2)\;dx + i\int_{-\infty}^\infty \sin(x^2)\;dx = (1+i)\sqrt\frac{\pi}{2}$$ • Thanks. I didn't know they had a name. I always call them 'gaussian integrals' :) – Tom Oct 3 '16 at 14:18 • This is the best solution; all of of the others, such commotion! Even myself, I would have expressed it as an error function, but this is much more direct. – Cye Waldman Jun 13 '17 at 18:37 The integral is a Fresnel one. Usually, it's evaluated in a 'pizza-slice' closed contour wich subtends a $\ds{\pi/4}$-angle in the complex plane first quadrant ( details can be seen elsewhere ). Since there are not any poles inside such contour, the integral along $\ds{\pars{0,\infty}}$ satisfies: \begin{align} \int_{0}^{\infty}\expo{\ic x^{2}}\dd x & = \lim_{r \to \infty}\bracks{-\int_{0}^{\pi/2} \exp\pars{\ic\bracks{r\expo{\ic\theta}}^{2}}r\expo{\ic\theta}\ic\,\dd\theta - \int_{r}^{0}\exp\pars{\ic\bracks{\xi\expo{\pi\ic/4}}^{2}}\expo{\pi\ic/4}\dd\xi} \\[5mm] & = -\ic\,\overbrace{\lim_{r \to \infty}\int_{0}^{\pi/2}\exp\pars{\ic r^{2}\cos\pars{2\theta}} \exp\pars{-r^{2}\sin\pars{2\theta}}r\expo{\ic\theta}\dd\theta}^{\ds{=\ 0}} \\[5mm] & \mbox{} + \expo{\pi\ic/4}\ \underbrace{\int_{0}^{\infty}\expo{-\xi^{2}}\,\dd\xi} _{\ds{=\ {\root{\pi} \over 2}}}\ =\ \bbox[8px,border:0.1em groove navy]{{\root{\pi} \over 2}\,\expo{\pi\ic/4}} \end{align} Note that \begin{align} 0 & < \verts{\int_{0}^{\pi/2}\exp\pars{\ic r^{2}\cos\pars{2\theta}} \exp\pars{-r^{2}\sin\pars{2\theta}}r\expo{\ic\theta}\dd\theta} < \int_{0}^{\pi/2}\exp\pars{-r^{2}\sin\pars{2\theta}}r\,\dd\theta \\[5mm] & = {1 \over 2}\,r\int_{0}^{\pi}\exp\pars{-r^{2}\sin\pars{\theta}}\dd\theta = r\int_{0}^{\pi/2}\exp\pars{-r^{2}\sin\pars{\theta}}\dd\theta < r\int_{0}^{\pi/2}\exp\pars{-\,{2r^{2} \over \pi}\,\theta}\dd\theta \\[5mm] & = {\pi \over 2}\,{1 - \expo{-r^{2}} \over r}\,\,\,\stackrel{r\ \to\ \infty}{\to} 0 \end{align} • you should use the complex analysis notation : $\int_0^\infty e^{i x^2}dx = \lim_{a \to 1,Im(a) > 0} \lim_{R \to \infty}\int_0^{a R}e^{i z^2}dz$ and $\int_0^{a R}+\int_{a R}^{e^{i \pi/4}R } +\int_{e^{i \pi/4}R}^0 e^{i z^2}dz = 0$ but $\lim_{R \to \infty}\int_{a R}^{e^{i \pi/4}R } e^{i z^2}dz = 0$ while $\lim_{R \to \infty}\int_{e^{i \pi/4}R}^0 e^{i z^2}dz = -e^{i \pi/4} \int_0^\infty e^{-x^2}dx$ – reuns Oct 9 '16 at 5:41 • Great contribution, thanks :) @Felix Marin – Tom Oct 9 '16 at 5:41 • @Tom It's nice it was useful for you. It's a standard evaluation you can check in Wikipedia ( I gave the link in the text ). Thanks. $\left(\bullet\,\,\,\,\, \bullet\atop {\mid \atop \smile}\right)$. – Felix Marin Oct 9 '16 at 5:45
2019-06-19T01:03:57
{ "domain": "stackexchange.com", "url": "https://math.stackexchange.com/questions/1952056/computing-the-integral-int-expix2-dx", "openwebmath_score": 0.9956888556480408, "openwebmath_perplexity": 421.2942365913725, "lm_name": "Qwen/Qwen-72B", "lm_label": "1. YES\n2. YES", "lm_q1_score": 0.9811668734137681, "lm_q2_score": 0.8539127529517043, "lm_q1q2_score": 0.8378309059817671 }
https://math.stackexchange.com/questions/2658728/for-which-values-of-x-does-these-series-converge-absolutely-converge-conditio
# For which values of $x$ does these series converge absolutely, converge conditionally or diverge I have these (math) problems where I am first supposed to find out if a given series converge or diverge, and then, by using that result, multiplying the same series by $x^n$ and find out for which values of $x$ the series converge absolutely, conditionally or diverge. I think I am starting to get a certain idea of which converge/divergence tests to use for different types of series, but the "by using that result"-part is confusing me a little bit. ### First series I concluded $\sum_{n=2}^{\infty}\frac{n^2 + 1}{n^2 - 1}$ diverges by applying the divergence test: $\lim_{n \to \infty}\frac{n^2 + 1}{n^2 - 1} = 1 \neq 0$. However, when I try to find out for which values of $x$ the series $\sum_{n=2}^{\infty}\frac{n^2 + 1}{n^2 - 1}x^n$ converge absolutely, conditionally or diverge, I dont know how to directly use the result I found previously. My book says the way to find out the radius of converge is by using the ratio test: $\frac{1}{R} = L = \lim_{n \to \infty}\left|\frac{a_{n + 1}}{a_n}\right| = \lim_{n \to \infty}\frac{((n+1)^2 +1)(n^2 - 1)}{((n+1)^2 - 1)(n^2 + 1)} = 1$ (by multiplying by $\frac{1}{n^2}$) So because radius of convergence is $1$, I concluded the series converges absolutely for $x \in (-1 , 1)$ and diverge for $(-\infty, -1]\cup[1, \infty)$. Here I knew already that the series diverges for $x = 1$ and thus can not converge absolutely for $x \leq -1$. It does also not converge conditionally as $\lim_{n \to \infty}\frac{n^2 + 1}{n^2 - 1}(-1)^n \neq 0$. So, in other words, assuming I am on the right track here, I need a little advice how to do this more direct if possible. I do one more example to see if I have thought about it in the right way.. ### Second series I used the comparison test to conclude that $\sum_{n=0}^{\infty}\frac{2^n}{4^n + 1}$ converges since $\frac{2^n}{4^n + 1} \lt \frac{1}{2^n}$ for $n = 0, 1, 2, \dots$ Now, trying to do the same as previously for $\sum_{n=0}^{\infty}\frac{2^n}{4^n + 1}x^n$, I don't see a way to use what I just found directly, so I do the ratio test again to find the radius of converge with center of converge at $x = 0$: $\frac{1}{R} = L = \lim_{n \to \infty}|\frac{a_{n + 1}}{a_n}| = \lim_{n \to \infty}\frac{2^{n + 1}(4^n + 1)}{2^n(4^{n + 1} + 1)} = \lim_{n \to \infty}\frac{2*4^n}{4*4^n + 1} + \lim_{n \to \infty}\frac{2}{4*4^n + 1} = \frac{2}{4} \Rightarrow R = 2$. So I concluded the series converges (absolutely) for $x \in (-2, 2)$ which actually makes sense when I look at it, as $x = 2$ would make the numerator $4^n$ and the series would diverge like $\sum_{n=0}^{\infty}1$. Again I do not see any way to make the series converge conditionally and thus diverges everywhere else but $x \in (-2, 2)$ ### Summary I need advice how I can use the information about whether a series converges or diverges to decide for which values of the same series multiplied by $x^n$ converges absolutely, converges conditionally or diverges. Any help appreciated, sorry if the math things are a little bit tiny. I don't think there's really a better way to do either of these. A slightly less direct but nicer way would be to use the limit comparison test to compare it with $$\sum_{n=2}^{\infty} x^n$$ (for the first series) or $$\sum_{n=2}^{\infty} (x/2)^n$$(for the second), but it's essentially the same idea. Generally speaking, if you can determine how your original series converges, you can find the radius convergence of it multiplied by $x^n$ in this manner. (Also, a note: If you put double-dollar-signs around expressions they center and become bigger - that's called displaymath mode and it's generally used for expressions with large operators.) • But if i use the comparison test, how do I know how the original series converges? My book says if I have series $\{a_n\}$ and $\{b_n\}$ then I know that $\{a_n\}$ converges if $L < \infty$ and $\{b_n\}$ converges or diverges if $L > 0$ and $\{b_n\}$ diverges where $L = \lim_{n \to \infty}\frac{a_n}{b_n}$. It says nothing about <b>how</b> – Amoz Feb 21 '18 at 12:02 • I am very curious because I am trying to do do the same thing for $$\sum_{n=1}^{\infty} \frac{n}{n^{2}\sqrt{n} + 1}$$ which wasn't so easy to do with the ratio test, but I know it behaves like $$\frac{1}{n^{\frac{3}{2}}}$$ – Amoz Feb 21 '18 at 16:50 • @Amoz The idea is that, if $L=\lim_{n\to\infty} \frac{a_n}{b_n}$ is a positive real number (i.e. $0<L<\infty$), then the series $\sum_{n=1}^{\infty} a_n$ converges iff the series $\sum_{n=1}^{\infty} b_n$ does. In other words, if their ratio tends to a nonzero finite constant, the two series behave the same way. It's usually used where one can determine whether $\sum_{n=1}^{\infty} b_n$ converges relatively simply. – Carl Schildkraut Feb 21 '18 at 21:27
2019-06-18T23:31:20
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https://math.stackexchange.com/questions/3757202/are-all-finite-dimensional-algebras-of-a-fixed-dimension-over-a-field-isomorphic/3758022
# Are all finite-dimensional algebras of a fixed dimension over a field isomorphic to one another? Suppose I have a finite-dimensional algebra $$V$$ of dimension $$n$$ over a field $$\mathbb{F}$$. Then $$V$$ is an $$n$$-dimensional vector space and comes equipped with a bilinear product $$\phi : V \times V \to V$$. Suppose now that I have another finite-dimensional algebra $$W$$ of dimension $$n$$ over $$\mathbb{F}$$ equipped with a bilinear product $$\psi: W \times W \to W$$. Certainly, $$V$$ and $$W$$ are isomorphic as vector spaces but are they isomorphic as $$\mathbb{F}$$-algebras? The question I'm really asking here is - Are all $$n$$-dimensional algebras over $$\mathbb{F}$$ isomorphic to one another? If the answer is yes, then this is my attempt at constructing such an isomorphism. Suppose I want to define an $$\mathbb{F}$$-algebra isomorphism between $$V$$ and $$W$$. To do this I'd need to define a map $$f : V \to W$$ such that • $$f(ax) = af(x)$$ for all $$a \in \mathbb{F}, x \in V$$ • $$f(x+y) = f(x) + f(y)$$ for all $$x, y \in V$$ • $$f(\phi(x, y)) = \psi(f(x), f(y))$$ for all $$x, y \in V$$ If $$\{v_1, \dots, v_n\}$$ and $$\{w_1, \dots, w_n\}$$ are bases for $$V$$ and $$W$$ respectively then both $$\phi$$ and $$\psi$$ being bilinear maps are completely determined by their action on basis vectors $$\phi(v_i, v_j)$$ and $$\psi(w_i, w_j)$$ for $$1 \leq i, j, \leq n$$. It turns out that $$\phi(v_i, v_j) = \sum_{k=1}^n \gamma_{i,j,k}v_k$$ and $$\psi(v_i, v_j) = \sum_{k=1}^n \xi_{i,j,k}w_k$$ for some collection of scalars $$\gamma_{i,j,k}$$ and $$\xi_{i,j,k}$$ called structure coefficients. So then if both the $$n^3$$ collections of scalars $$\gamma_{i,j,k}$$ and $$\xi_{i,j,k}$$ are all non-zero then we can define $$f : V \to W$$ by $$f(a_1v_1 + \cdots + a_nv_n) = a_1 \frac{\xi_{i,j,1}}{\gamma_{i,j,1}}w_1 + \cdots + \frac{\xi_{i,j,n}}{\gamma_{i,j,n}}w_n$$ and it will turn out that $$f$$ is the desired isomorphism of algebras as one can then check that $$f(\phi(v_i, v_j)) = \psi(w_i, w_j) = \psi(f(v_i), f(v_j))$$ for all $$i$$ and $$j$$. However what if it's the case that for $$\phi$$ some $$\gamma_{i, j, k}$$ is zero and the corresponding $$\xi_{i, j, k}$$ is non-zero? I don't see any way to get an isomorphism in that case. Is it still possible to construct an isomorphism in that case? • The right side of the equation has $i,j$, not summed or quantified, while the left side has no $i,j$. What do you mean by this? – mr_e_man Jul 15 '20 at 14:15 • Everyrone is working way to hard to provide counter examples. What about the field $\mathbb{F}$ vs. the additive group of $\mathbb{F}$ with $0$-multiplication? – Jason DeVito Jul 15 '20 at 15:59 • @JasonDeVito Giving non-unital counterexamples is cheating! – Earthliŋ Jul 15 '20 at 16:33 • @Earthliŋ: Perhaps it is, but I made sure to read the question for the word "unit" before posting my comment. (I also purposefully didn't provide it as an answer.) – Jason DeVito Jul 15 '20 at 16:34 • $\Bbb{Q}(\sqrt2)$ and $\Bbb{Q}(\sqrt3)$ are not isomorphic as $\Bbb{Q}$-algebras in spite of both being 2-dimensional. – Jyrki Lahtonen Jul 16 '20 at 21:36 They will not necessarily be isomorphic. Consider $$V = \mathbb F[x] / (x^n)$$ and $$W = \mathbb F^n$$ with componentwise multiplication.These are both $$n$$ dimensional $$\mathbb F$$ algebras. However, $$V$$ contains a nilpotent element, $$x$$, whereas $$W$$ contains no nilpotent elements. Indeed, if we had an $$\mathbb F$$-algebra homomorphism $$f: V \longrightarrow W$$ then as $$0 = f(x^n) = f(x)^n$$, we'd need $$f(x) = 0$$ so any map between the two must have a nontrivial kernel. Another very familiar example: $$\mathbb{C}\neq\mathbb{R}\times \mathbb{R}$$. The complex numbers are a field, but $$(1,0)(0,1)=(0,0)$$ in $$\mathbb{R}\times \mathbb{R}$$, so it has non-trivial zero-divisors. In general, the answer is "no", even if one requires $$V$$ and $$W$$ to be fields. For example, the rings $$\mathbb{Q}(\sqrt{2})$$ and $$\mathbb{Q}(\sqrt{3})$$ are two non-isomorphic fields that both have dimension $$2$$ over $$\mathbb{Q}$$. • I guess if you require $V$ and $W$ both to be fields, then this is true iff $V$ is algebraically closed, real closed or pseudofinite. – tomasz Jul 15 '20 at 0:51 There are several answers that point out why the statement of the question cannot be true, probably the simplest example being $$\Bbbk [x] / (x^2) \not\simeq \Bbbk \times \Bbbk$$. Classifying all finite-dimensional algebras of a given dimension is actually rather involved and very far from being just one algebra in each dimension. Note that you can even come up with finite-dimensional noncommutative algebras. For example, from the quiver $$\bullet \to \bullet$$ you can build a noncommutative $$3$$-dimensional algebra with $$\Bbbk$$-basis $$e_1, e_2, \alpha$$, where • $$e_1, e_2$$ are viewed as "constant paths" at the vertices, which are orthogonal idempotents, i.e. $$e_i e_j = \delta_{ij}$$ • $$\alpha$$ is viewed as corresponding to the arrow and $$e_1, e_2$$ are viewed as "identities at" the vertices, so $$e_1 \alpha = \alpha$$ and $$\alpha e_2 = \alpha$$ • the product of paths which cannot be composed are defined to be $$0$$ in this algebra, so $$e_2 \alpha = \alpha e_1 = \alpha^2 = 0$$ and extending these rules linearly gives a well-defined associative multiplication. More generally, you can take the path algebra of any quiver and quotient by any two-sided ideal, which if you choose the ideal correctly will give a finite-dimensional algebra, which is usually non-commutative. Finite-dimensional algebras can be studied via their categories of finite-dimensional modules (which in some cases can actually be described rather explicitly) and it turns out that the construction of finite-dimensional algebras via quivers gives all algebras up to Morita equivalence (i.e. using quivers you find the module categories of all finite-dimensional algebras). Let $$G$$ and $$H$$ be finite groups of the same order, such that $$G$$ is abelian and $$H$$ is not. Then the group rings $$V = \mathbb{F}G$$ and $$W = \mathbb{F}H$$ share the same dimension, but $$V$$ is commutative while $$W$$ is not. https://en.wikipedia.org/wiki/Group_ring No. For example, $$\mathbf Q[\sqrt n]$$ are pairwise nonisomorphic (where $$n$$ ranges over squarefree integers distinct from $$1$$), but all have dimension $$2$$ over $$\mathbf Q$$. In general, if $$K$$ is not algebraically closed, then it admits a finite algebraic extension $$L\supsetneq K$$, and then $$L$$ and $$K^{[L:K]}$$ have the same dimension and are not isomorphic. Even if $$K$$ is algebraically closed, $$K^4$$, $$M_{2\times 2}(K)$$ and $$K[x]/(x^4)$$ are non-isomorphic four-dimensional algebras over $$K$$. Edit: as suggested in the comments, the counterexamples listed above are essentially all the counterexamples. More precisely, the answer is yes if you restrict yourself to finite-dimensional algebras over an algebraically closed field $$k$$ which are reduced (contain no nilpotents). In other words, the only such algebras are of the form $$k^n$$. One can show that in this case, the algebra $$A$$ is semisimple (because it is Artinian and the Jacobson radical is zero), so by Wedderburn's theorem, it follows that it is a product of matrix rings over division algebras. Since there are no finite-dimensional division algebras over $$k$$ (because $$k$$ is algebraically closed), and no proper matrix ring is reduced (because it contains strictly upper triangular matrices, which are nilpotent), it follows that $$A\cong k^n$$ for some $$n$$. • $\mathbb{C}[x]/\{x^2\}$ is finite dimensional and commutative, but not isomorphic to $\mathbb{C}\times \mathbb{C}$. – tkf Jul 15 '20 at 0:15 • @tkf: You're right. For some reason, I thought that such an algebra would automatically be semisimple, but your example shows that this is not the case. I removed the erroneous claim. – tomasz Jul 15 '20 at 0:18 • I think you just need to add the condition that it does not contain nilpotent elements. – tkf Jul 15 '20 at 0:19 • @tkf: I guess commutative+reduced is sufficient, because that makes the algebra the ring of regular functions of a zero-dimensional variety. But I don't see whether or not reduced finite-dimensional implies commutative... – tomasz Jul 15 '20 at 0:35 • @tkf: I meant over an algebraically closed field. – tomasz Jul 15 '20 at 15:34
2021-06-13T08:55:20
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https://ckrao.wordpress.com/2012/04/10/integers-equal-to-the-sum-of-three-cubes/
# Chaitanya's Random Pages ## April 10, 2012 ### Integers equal to the sum of three cubes Filed under: mathematics — ckrao @ 11:32 am One easy-to-state open problem in number theory is whether every integer not of the form $9k \pm 4$ is the sum of three cubes. Some solutions are easy to find, while others are notoriously difficult. For example 29 can be written as $3^3 + 1^3 + 1^3$, or even $4^3 - 3^3 - 2^3$. It is clear that $9k \pm 4$ cannot be written as the sum of three cubes since $x^3 \equiv -1, 0 \text{ or } 1 \mod 9$. There also exist infinite families of solutions in certain cases (though these need not represent all solutions). For example, $\displaystyle (9t^3 + 1)^3 + (9t^4)^3 + (-9t^4 - 3t)^3 = 1$ $\displaystyle (6t^3 + 1)^3 - (6t^3 - 1)^3 - (6t^2)^3 = 2$ For a long time it was not known whether 30 can be written as the sum of cubes. The following is the “simplest” way this can be done, found as recently as 1999. $\begin{array}{lcl} 30 &=& (2,220,422,932)^3 + (-2,218,888,517)^3 + (-283,059,965)^3\\ &=& 10,947,302,325,566,084,787,191,541,568 - 10,924,622,727,902,378,924,946,084,413 - 22,679,597,663,705,862,245,457,125 \end{array}$ It is still not known whether 33, 42 or 74 can be written as the sum of three cubes (edit: it is now known that 74 can be – see [5])! In these cases no solution has been found when any of the cubes is less than $(10^{14})^3$ in magnitude [4]. One of the computational algorithms to find solutions to $a^3 + b^3 + c^3 = n$ where $n$ is relatively small, proposed by Elkies, involves converting the equation to $(-a/c)^3 + (-b/c)^3 \approx 1$ and then considering rational points close to the curve $y = \sqrt[3]{1-x^3}, x \in [0, 1/\sqrt[3]{2}]$. This curve is covered by small parallelograms and the problem is converted to finding lattice points in a pyramid using basis reduction followed by the Fincke-Pohst algorithm [4]. #### References [2] Hisanori Mishima, Chapter 4. n=x^3+y^3+z^3 [3] D.J. Bernstein, http://cr.yp.to/threecubes.html [4] A.-S. Elsenhans and J. Jahnel, New sums of three cubes, Math. Comp. 78 (2009), 1227–1230, available here. [5] Sander G. Huisman, Newer sums of three cubes, http://arxiv.org/abs/1604.07746, April 2016. 1. Hmm, wondering what the name for this equation is: a^3+b^3+c^3=z^3 Or the differences between the differences between consecutive cubes. Or for that matter the proof that if the differences process is recursively used on any give power of x (y=x^n, x is within positive integers), that the result is always factorial based on n. In other words, it equals ‘n!’. This HAS to have been looked into during the days when machines designed to do that automatically were being used around the time after Colin Maclaurin died. Comment by Joe — November 1, 2013 @ 12:07 pm • Not sure, but it is referred to a 3.1.3 equation here: http://mathworld.wolfram.com/DiophantineEquation3rdPowers.html To understand why the n! result is true, refer to the final equation of this earlier blog post and let $p(x) = x^n$. Match the coefficients of $x^n$ of both sides of that equation to conclude that $a_n = n!$ where $a_n$ is the bottom number of the finite difference table. Comment by ckrao — November 2, 2013 @ 11:24 am • Well I kind of figured it was related to the n’th derivative function for any given n’th power of x. Forgive my attempt at doing this in text: Dx^5 of x^5=5*4*3*2*1 This is true since Dx x^5=5x^4 and you can repeat the process for all powers until you get x^0 (in other words, a constant). What would be the name for a generalized equation family like this? a_1^n + … + a_n^n = b^n (where _ means subscript which is a way of representing an array member) So sum of 4 hypercubes having a ‘content’ of exactly 1 hypercube. That is the geometric description of a^4+b^4+c^4+d^4=z^4 obviously. Change n to 2 and you get good ol’ Pythagoras’ Theorem. Which of course has a geometric description exactly matching a circle (constant radius) or triangle (hypotenuse changes). 😉 Thanks for answering earlier. Diophantine equations are… a bit broad for something that specific. It’s sort of the other side of Fermat’s (big) Last Theorem. It could be fun (or frustrating – haven’t checked if anyone’s succeeded for this specific case) to prove there’s infinite solutions for some families of this nature. Many arbitrary polynomials have non-algebraic solutions… Comment by Joe — November 2, 2013 @ 11:26 pm Comment by Joe — November 2, 2013 @ 11:33 pm 2. sum of three different cube of positive integers is equal to a another cube of positive integers Comment by Subhankar Mondal — February 28, 2014 @ 7:19 am 3. 74 = 66229832190556³ + 283450105697727³ – 284650292555885³ Comment by John Doe — June 1, 2016 @ 12:35 pm • Many thanks for this! The post has now been updated with the link you have provided. Comment by ckrao — June 19, 2016 @ 5:39 am 4. 21^3 + 28^3 + 35^3 = 42^3 , which is the series 3 4 5 6 multiplied by 7. by a program I wrote. There are so many sums of cubes; perhaps they follow a rule similar to the Pythagorean theorem! (multiples work) Comment by — August 2, 2016 @ 8:38 pm • But we talk about $a^3+b^3+c^3=42$, not about $a^3+b^3+c^3=42^3$. Comment by guest — July 8, 2018 @ 5:07 am 5. and k stands for… Comment by mos — August 19, 2018 @ 3:06 am • an integer Comment by ckrao — September 15, 2018 @ 6:33 am 6. Could somebody please explain in further detail why the values that can be represented as 9k+-4 can’t be solved. I’m struggling to wrap my head around why. Comment by sunny — October 15, 2018 @ 1:51 pm • Yes, a perfect cube has remainder -1, 0 or 1 when divided by 9 so the sum of three of these cannot be of the form 9k+-4. Comment by ckrao — October 16, 2018 @ 1:35 am Blog at WordPress.com.
2019-01-20T18:23:08
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https://math.stackexchange.com/questions/238752/question-on-probability-using-bayes-theorem/238770
# Question on probability using Bayes' theorem In a certain day care class, $30\%$ of the children have grey eyes, $50\%$ of them have blue and the other $20\%$'s eyes are in other colors. One day they play a game together. In the first run, $65\%$ of the grey eye ones, $82\%$ of the blue eyed ones and $50\%$ of the children with other eye color were selected. Now, if a child is selected randomly from the class, and we know that he/she was not in the first game, what is the probability that the child has blue eyes? My solution Let's say $B =$ blue, $G =$ grey and $O =$ "Other color" and $NR =$ "not selected for the first run" $$P(B \mid NR) = \frac{P(NR \mid B)P(B)}{P(G)P(NR \mid G) + P(B)P(NR \mid B) + P(O)P(NR \mid O)}$$ On substituting values $$P(B \mid NR) = \frac{0.5 \cdot (1-0.82)}{(0.3 \cdot (1-0.65)) + (0.5 \cdot (1-0.82)) + (0.2 \cdot (1-0.5))}$$ $$P(B \mid NR) = 0.305$$ Is this the right way to use bayes theorem? Yes, this is correct. The general form of Bayes' rule is $$P(A_i|B)= \frac{P(B|A_i)P(A_i)}{\sum_jP(B|A_j)P(A_j)}$$ as you've used above. • Thanks for the generalized version of Bayes' theorem. Do you see anything wrong with the way I have arrived at the individual probabilities? Nov 16, 2012 at 17:12 • Nothing wrong that I can see... Nov 16, 2012 at 18:49 • Has that helped, do you have the answer now? Nov 16, 2012 at 19:23 Suppose 100x total students Blue eyed =50x Grey eyed =30x Other eyed =20x 1st run selection : Blue eyed =41x Grey eyed =19.5x Other eyed=10x Not selected in 1st round Blue= 9x Grey=10.5x Other =10x P=(9x)/(9x+10.5x+10x)=.305 $$p(\textrm{not game 1} \mid \textrm{blue}) \div p(\textrm{not game 1}) = 0.18×0.50 \div (0.18×0.50+0.5×0.2+0.35×0.30) =0.3051$$
2022-05-21T13:46:40
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https://math.stackexchange.com/questions/1574622/symmetric-transitive-sets
# Symmetric & Transitive Sets Let $A={a,b,c,d,e,f}$ and let $R\subset A\times A$ be a relation which is symmetric and transitive. You have been given some partial information about the relation which is that the following are known to be true for the relation: $$R(e,a), R(b,f), R(e,c), R(d,a), R(b,a)$$ Given all the above information does $R(c,f)$ hold or not? Explain your answer using the information you have been provided with above. I would argue that $R(c,f)$ does not hold. No two pairs would equal $R(c,f)$ and $R(f,c)$. Does this seem correct? Update: Just thought of this. R(b,f) and R(b,a) would be the same as R(f,a). R(f,a) and R(e,a) would be the same as R(f,e). Therefore, R(f,e) and R(e,c) would be the same as R(f,c). As the relation is symmetric, R(f,c) -> R(c,f). How's that? • That's correct! Note - be careful not to mix up symmetry and transitivity; on the last line, you meant to say symmetry. Dec 14, 2015 at 3:48 • Just corrected that! Did I solve this as efficiently as possible or is there a better method? – M-R Dec 14, 2015 at 3:50 • There is a technical term "transitive set", and while it relates to your use here, what you really mean is a transitive (and symmetric) relation. Dec 14, 2015 at 4:34 • I think if you draw a picture, you'll see that you've taken the most efficient path from $c$ to $f$, so in that sense this is the most efficient route using the properties of transitivity and symmetry. In another sense, you could argue it would be more 'neat' to simply draw the picture, and say "there is a path from $c$ to $f$, therefore $R(c, f)$." Efficiency isn't really what you're looking for in mathematics answers, it's more about understanding and effective communication. If you can explain your answer well, you've succeeded. Dec 14, 2015 at 20:46 This is a great question to apply some methods from Polya's 'How to Solve it'. Namely, can you draw a picture??? Because the relation $R$ is symmetric and transitive, it makes a lot of sense to draw this like a mathematical graph (that's different from a 'graph' of $f(x)=x^2$, by the way - wiki it). So if you draw $a, b, c, d, e, f$ as points on your page, and connect them with lines whenever you know they're related, the question becomes - is there a path from $c$ to $f$? If there was any such path, using, for example, $c\mapsto e\mapsto a\mapsto b\mapsto f$, then you would know, using symmetry: $$R(c, e), R(e, a), R(a, b), R(b, f)$$ Applying transitivity to the first pair, you could 'walk' from $c$ to $a$, concluding $R(c, a)$. Then repeat, until you conclude $R(c, f)$! Each of the following must hold. I leave as an exercise finding the justifications, in order. 1. $R(c,e)$ 2. $R(c,a)$ 3. $R(a,b)$ 4. $R(c,b)$ 5. $R(c,f)$
2022-05-28T10:53:34
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https://math.stackexchange.com/questions/2030634/how-can-we-factorize-x4-2x2-49-with-coefficients-in-mathbbr
How can we factorize $x^4 - 2x^2 + 49$ with coefficients in $\mathbb{R}$? [duplicate] How can we factorize ${x^4} - 2{x^2} + 49$ with coefficients in $\mathbb{R}$? A problem would be easier if it was a quadratic equation - we could simply find the roots and get the linear factors. Moreover, the polynomial $x^4 - 2x^2 + 49$ does not have a real root which would be easy to guess. (If we have one root, we could divide by linear factor determined by this root.) WolframAlpha says that this can be factorized as $$x^4 - 2x^2 + 49 = (x^2-4x+7)(x^2+4x+7).$$ But how can we get to this factorization? Factorization of this polynomial also appears as an example in an answer tp another question: Does the Rational Root Theorem ever guarantee that a polynomial is irreducible? marked as duplicate by Namaste, Stefan Mesken, Bill Dubuque, астон вілла олоф мэллбэрг, R_DNov 26 '16 at 5:40 • @costrom - If ${x^4} - 2{x^2} + 49 = 0$ then, this equation is not the answer in $\mathbb{R}$. – Under sky Nov 25 '16 at 20:12 The idea is try to complete square using the terms $x^4$ and $49$ so let's reorganize the expression: $x^4+14x^2+49-16x^2=(x^2+7)^2-(4x)^2=(x^2-4x+7)(x^2+4x+7)$ • - Thanks. ......... – Under sky Nov 25 '16 at 20:23 • If the downvoters could please tell me why I deserve that I could improve this and the futures anwers. – Arnaldo Jan 8 '17 at 11:24 • I think we may all have been downvoted without comment. :-/ – Cameron Buie Jan 8 '17 at 15:04 • @Cameron Buie: I also think so. Sadly, that behavior doesn't help anybody. Neither we nor the OP or even the downvoter. – Arnaldo Jan 8 '17 at 15:20 There are a few possible forms for such a factorization: • as a product of four (not necessarily distinct) linear factors with real coefficients, • as a product of two quadratic factors with real coefficients, neither of which has real roots, • as a product of two linear factors and a quadratic factor, all with real coefficients, such that the quadratic factor has no real roots. Can you see why there are no other possibilities? Now, note that $$x^4-2x^2+49=(x^2)^2-2x^2+1+48=(x^2-1)^2+48\ge 48$$ for all real $x,$ so that the polynomial has no real roots. This rules out the possibility of any linear factors with real coefficients, so the only possibility remaining is that we can factor it as a product of quadratics with real coefficients. Set $$x^4-2x^2+49=(px^2+qx+r)(sx^2+tx+u)$$ with $p,s$ non-zero. We must have $ps=1.$ (Why?) Without loss of generality, we can assume that $p=s=1,$ because if not, we can factor out $p$ from the first quadratic and $s$ from the second to get $$x^4-2x^2+49=\left(x^2+\frac qpx+\frac rp\right)\left(x^2+\frac tsx+\frac us\right),$$ which is of the form $$x^4-2x^2+49=(x^2+q'x+r')(x^2+t'x+u').$$ Now, starting from $$x^4-2x^2+49=(x^2+qx+r)(x^2+tx+u),$$ expand the right-hand side and gather like terms. This will yield a system of $4$ equations in $4$ variables, which you should try to solve. If it has real solutions, we've succeeded! If not, it isn't possible. With $t=x^2$, $$t^2-2t+49$$ has the roots $$t=1\pm i\sqrt{48}$$ and $$x=\pm\sqrt{1\pm i\sqrt{48}}=\pm a\pm ib.$$ Then $$(x+a+ib)(x-a+ib)(x+a-ib)(x-a-ib)=\\ (x^2+2ax+a^2+b^2)(x^2-2ax+a^2+b^2)$$ $$x^4-2x^2-49 = x^4+14x^2+49-16x^2$$ =$$(x^2+7)^2-16x^2$$ =$$(x^2-4x+7)(x^2+4x+7)$$. And we are done
2019-07-21T10:30:21
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https://math.stackexchange.com/questions/1493156/fixed-point-iteration-method-starting-point
# Fixed Point Iteration Method - Starting Point If $g(x):=21^{1/2}x^{-1/2}$, then $21^{1/3}$ is a fixed point of $g$. Question: Using the Fixed Point Iteration Method, are there conditions on the starting point $x_0$ in order for the method to converge? Justify. So it seems like any $x_0>0$ should be such that we have convergence. However, how to justify it? Geometrically, this seems plausible because of the curvature of $g$. Is there a criterion based on the second derivative of $g$? In theory one could solve the recurrence relation $$x_{n+1}=21^{1/2}x_n^{-1/2}$$ and find, according to WolframAlpha, $$x_n=21^{\frac{1}{3}(1-(-\frac{1}{2})^n)}e^{c(-(-1)^n)2^{1-n}}$$ where $c$ is a parameter depending on $x_0$, so that $x_n\to21^{1/3}$ no matter what $x_0>0$ we choose. However, in practice, I have no idea how to solve this type of recurrence relation... ## 1 Answer This is a bit problem-specific, but this trick appears to work. Suppose $x_0=a 21^{1/3}$ for some $a>0$. Then $x_1=21^{1/2} a^{-1/2} 21^{-1/6}=a^{-1/2} 21^{1/3}$. Can you now solve the recurrence by induction? What does your solution tell you? As for trying to do something more general, I think you might be able to show that if $f$ is $C^1$, decreasing, convex, has a fixed point, and its derivative is bigger than $-1$ at its fixed point, then the fixed point iteration will spiral inward toward the root. The idea of that would be to show that the even and odd subsequences are bounded and monotone. • Yes I think I can solve the recurrence by induction now: $x_n=a^{(-\frac{1}{2})^n}21^{1/3}$ so that $x_n\to a^0 21^{1/3}=21^{1/3}$. Since $a>0$ was arbitrary, so was $x_0>0$. Very clever. – Guest Oct 23 '15 at 1:14
2019-10-14T03:49:21
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https://www.hpmuseum.org/forum/thread-15683.html
(41C) Method of Successive Substitutions 10-04-2020, 04:21 PM Post: #1 Eddie W. Shore Senior Member Posts: 1,225 Joined: Dec 2013 (41C) Method of Successive Substitutions Repeating the Calculation Again and Again With an aid of a scientific calculator, we can solve certain problems of the form: f(x) = x. Examples include: tan cos sin x = x e^-x = x atan √x = x sin cos x = x (acos x) ^ (1/3) = x In any case with trigonometric functions, the angle mode will need to be in radians. You will also need a good guess to get to a solution and to know that at some real number x, f(x) and x intersect. Take the equation ln(3*x) = x with initial guess x0 = 1.512 Depending on the operating of the scientific calculator the keystrokes would be: AOS: 1.512 [ = ] Loop: [ × ] 3 [ = ] [ ln ] RPN: 1.512 [ENTER] Loop: 3 [ × ] [ ln ] ALG: 1.512 [ENTER/=] Loop: ln( 3 * Ans) [ ENTER/= ] Repeat the loop as many times as you like and hope you start seeing the answers converge. After repeating the loop over and over and over again, at six decimal answers, the readout will be about 1.512135. An approximate answer to ln(3x) = x, x ≈ 1.512134552 If your calculator has a solve function, you can check the answer, but this method can be useful if your calculator does not have a solve function. The program SUCCESS illustrates this method. HP 41C/DM41 Program: SUCCESS This program calls on the subroutine, FX. FX is where you enter f(x). End FX with the RTN command. Code: 01 LBL^T SUCCESS 02 ^T F<X>=X 03 AVIEW 04 PSE 05 ^T GUESS? 06 PROMPT 07 STO 00 08 ^T PRECISION? 09 PROMPT 10 STO 02 11 0 12 STO 03 13 1 14 STO 04 15 LBL 01 16 RCL 00 17 XEQ ^FX 18 STO 01 19 RCL 00 20 - 21 ABS 22 STO 04 23 RCL 01 24 STO 00 25 1 26 ST+ 03 27 200 28 RCL 03 29 X>Y? 30 GTO 02 31 RCL 02 32 CHS 33 10↑X 34 RCL 04 35 X>Y? 36 GTO 01 37 ^T SOL= 38 ARCL 01 39 AVIEW 40 STOP 41 ^T ITER= 42 ARCL 03 43 AVIEW 44 STOP 45 ^T DIFF= 46 ARCL 04 47 AVIEW 48 STOP 49 GTO 04 50 LBL 02 51 ^T NO SOL FOUND 52 AVIEW 53 STOP  54 LBL 04 55 END Examples for FX: f(x) = sin cos x. Program: LBL ^FX COS SIN RTN f(x) = e^-x Program: LBL ^FX CHS E↑X RTN Be aware, some equations cannot be solved in this manner, such as x = π / sin x and x = ln(1 / x^4). Cheung, Y.L. "Using Scientific Calculators to Demonstrate the Method of Successive Substitutions" The Mathematics Teacher. National Council of Teachers of Mathematics. January 1986, Vol. 79 No. 1 pp. 15-17 http://www.jstor.com/stable/27964746 Blog entry: https://edspi31415.blogspot.com/2020/10/...od-of.html 10-06-2020, 12:57 AM (This post was last modified: 10-06-2020 02:05 AM by Albert Chan.) Post: #2 Albert Chan Senior Member Posts: 1,601 Joined: Jul 2018 RE: (41C) Method of Successive Substitutions (10-04-2020 04:21 PM)Eddie W. Shore Wrote:  Take the equation ln(3*x) = x with initial guess x0 = 1.512 Sometimes, convergence may be slow, or not at all. We can place a weight on it. With my Casio FX-115MS 1.512 = ln(3 Ans  =       → 1.512045566 =       → 1.512075703 =       → 1.512095633 r = (95633-75703) / (75703-45566) = 19930 / 30317 Convegence is slow (we wanted small |r|) Assume same trend continued (constant r), estimated converged to: 1.512045566 + 0.000030317/(1-r) = 1.512134548 Let's check if assumption is good. Continued on ... w = 1/(1-r) ≈ 3       x = (1-w)*x + w*ln(3*x) = -2 x + 3 ln(3*x) -2 Ans + 3 ln( 3 Ans  =       → 1.512135175 =       → 1.512134542 =       → 1.512134552, converged
2021-10-26T21:49:27
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https://grindskills.com/what-are-the-minimum-and-maximum-values-of-variance-closed/
# What are the minimum and maximum values of variance? [closed] I am new to statistics. I am getting my hands dirty on VarianceThreshold. I am having a single dimensional array, containing N values. What’s the maximum and minimum values of a variance for any values present in array? I guess the minimum value will always be non-negative but I don’t know about the maximum value. I have googled it but couldn’t find a good answer. I interpret this question as asking Given a set of $$N$$ numbers $$x_1, x_2, \ldots,x_N$$, what is the minimum and maximum values that the variance $$V$$, defined as $$V = \frac 1N \sum_{k=1}^N (x_k-\bar{x})^2 ~~ \text{where}~\bar{x}=\frac 1N \sum_{k=1}^N x_k$$ can take on? Well, the minimum value of $$V$$ is $$0$$ as Daniel Lopez’s comment points out, and it occurs if and only if all the $$N$$ numbers have the same value. At the other end, every finite set of real numbers has a (finite) upper bound (call it $$b$$) and a (finite) lower bound (call it $$a$$), and $$V \leq \frac{(b-a)^2}{4} = \left(\frac{b-a}{2}\right)^2 = \left(\frac{\mathcal R}{2}\right)^2\tag{1}$$ where $$\mathcal R$$ is the range of the set of $$N$$ numbers. Note that it is not necessary to know the values of $$b$$ and $$a$$ separately; we only need the range $$\mathcal R = (b-a)$$ to calculate the upper bound $$(1)$$ on $$V$$. If $$N$$ is an even number, there exist sets for which the bound $$(1)$$ holds with equality: these are sets for which $$\frac N2$$ of the $$x_k$$ have value $$b$$ and the other $$\frac N2$$ have value $$a$$. For odd $$N$$, the bound $$\frac{(b-a)^2}{4}$$ still applies but cannot be attained with equality if $$N>1$$. For odd $$N>1$$, the maximum value is $$\frac{(b-a)^2}{4}\frac{N^2-1}{N^2}$$ and is attained by a set in which $$\frac{N-1}{2}$$ of the $$x_k$$ have value $$b$$, $$\frac{N+1}{2}$$ have value $$a$$, or vice versa. For details, see this answer of mine on math.SE. In another answer (and the comments on it), @Jim has argued that “A set of $$N$$ real numbers” tells the listener nothing whatsoever about the set if you don’t know what any of the values are, even the minimum or maximum, and so the only completely correct answer is that maximum possible $$V$$ is unbounded: any other answer (such as mine above or a couple of possible answers suggested by Jim) must be be festooned with caveats that the answer is based on assumptions might be unwarranted. I disagree. Even if a secretive questioner is unwilling to share any details about the set he/she is concerned about, my answer gives the questioner enough information to find the maximum possible value of $$V$$ for him/herself from very minimal information abut the set: just the range suffices, no need to know even what $$N$$ is! EDIT: (by AHK) Corrected the maximum variance for odd $$N>1$$ and the corresponding choice of $$x_i$$‘s
2022-09-29T02:04:54
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https://math.stackexchange.com/questions/2842674/computing-double-integral-using-linear-algebra
# Computing double integral using linear algebra $$\iint_D (6x+2y) \, \mathrm d x \,\mathrm d y$$ where $$D$$ is the convex hull of $$4$$ given points, $$D = \mbox{conv} \left\{ (0,0),(-2,6),(3,2),(1,8) \right\}$$ This is a parallelogram with "unit vectors" $$(-2,6)$$, $$(3,2)$$. I wanted to give a try to solve the following problem with algebra instead of calculus. So, I thought about calculating the area of it and piling it up to get its volume. I got the cosine between both vectors with the dot product formula, which is $$\dfrac{6}{\sqrt{13\cdot40}}$$ and the sine with the Pythagoras identity which is $$\sqrt{\dfrac{484}{13\cdot40}}$$ I remember that $$|a|\cdot |b|\cdot \sin\alpha$$ gives the height, and $$|a|\cdot |b|\cdot \cos\alpha$$ gives the area. So I figured that maybe this would solve the integral problem for the area? $$\dfrac{6}{\sqrt{13\cdot40}} \cdot \sqrt{13} \cdot \sqrt{40}\sqrt{\dfrac{484}{{13\cdot40}}}\cdot \sqrt{13} \cdot \sqrt{40}$$ That gives me $$6 \cdot 22 = 132$$ That's wrong but the right result is $$11 \cdot 22 = 242$$ so there might be something in that? • You can compute an integral with the aid of areas if your integrand is 1 or a constant. In your case you have an integrand which contains other variables $x,y$... In this case you should decompose your domain and transform the double integral into two successive one dimensional integrals. Or change variables in order to transform your domain into a square (not sure that's easier to do...) Jul 6, 2018 at 9:37 • May be the solid volume of which you want to get is not a parallelipiped. Its top is not parallel to XOY Jul 6, 2018 at 9:40 • @BeniBogosel: would it work if I did a change of base from \begin{pmatrix} \begin{align} 3 & -2 \\ 2 && 6\\ \end{align} \end{pmatrix} Sorry, I can't separate the 2 and the 6 properly! Jul 6, 2018 at 10:03 • Maybe if I made a cube of the integral then I would be able to calculate it's volume then? Jul 6, 2018 at 10:12 The area of the domain is given by the cross product of two sides, $$(2,-6)\times(3,2)=22.$$ Then the integral must be the area times the average value of the integrand, which by linearity is also the value of the function at the centroid, $\left(\dfrac12,4\right)$, hence $$22\left(6\cdot\frac12+2\cdot4\right)=22\cdot11.$$ Your answer cannot be right, because you make no use of the coefficients $6$ and $2$ in the integrand. What you are evaluating is $$|a|^2|b|^2\cos\alpha\sin\alpha$$ which has no justification. Another solution is by noting that the volume is a truncated prism with bases $(0,0,0),$ $(-2,6,0),$ $(3,2,0),(1,8,0)$ and $(0,0,0),(-2,6,0),(3,2,22),(1,8,22)$, which is the half of a non-truncated prism of height $22$. Hence $$22\cdot\frac{22}2.$$ • Thanks @YvesDaoust, I really love your answer! How do I find the centroid? Jul 6, 2018 at 10:46 • @Dovendyr: by symmetry, the centroid of the parallelogram is the centroid of the four vertices, i.e. the arithmetic mean. – user65203 Jul 6, 2018 at 10:46 • I like your second take on it as well, but can you send some graphic representation?? Jul 6, 2018 at 10:47 – user65203 Jul 6, 2018 at 10:47 • I got an empty box. Even though I follow the instructions! Graphics3D[ Prism[ {0, 0, 0}, {-2, 6, 0}, {3, 2, 0}, {1, 8, 0}, {0, 0, 0}, {-2, 6, 0}, {3, 2, 22}, {1, 8, 22}]] Jul 6, 2018 at 10:56 Let $$\begin{bmatrix} x\\ y\end{bmatrix} = \begin{bmatrix} 3 & -2\\ 2 & 6\end{bmatrix} \begin{bmatrix} u\\ v\end{bmatrix}$$ where the determinant of the matrix is $22$. Since $$6x + 2y = \begin{bmatrix} 6\\ 2\end{bmatrix}^\top \begin{bmatrix} x\\ y\end{bmatrix} = \begin{bmatrix} 6\\ 2\end{bmatrix}^\top \begin{bmatrix} 3 & -2\\ 2 & 6\end{bmatrix} \begin{bmatrix} u\\ v\end{bmatrix} = \begin{bmatrix} 22\\ 0\end{bmatrix}^\top \begin{bmatrix} u\\ v\end{bmatrix} = 22 u$$ we have $$\iint_D (6x+2y) \, \mathrm d x \,\mathrm d y = \iint_{[0,1]^2 } (22 u) (22 \, \mathrm d u \,\mathrm d v) = 22^2 \int_0^1 u \, \mathrm d u = \frac{22^2}{2} = 22 \cdot 11$$ • This is amazing! So neat! I have to give it some thought to understand how and why the determinant comes into play. I'll be back! Jul 6, 2018 at 13:15 • Note that the gradient of $6x+2y$ is orthogonal to the 2nd column of the matrix. The integrand was not chosen randomly! Jul 6, 2018 at 13:35 • I have been thinking a bit about that, but why is that : $$\begin{bmatrix} u\\ v\end{bmatrix} = \begin{bmatrix} 3 & -2\\ 2 & 6\end{bmatrix} \begin{bmatrix} x\\ y\end{bmatrix}$$? Isn't it the vectors $u$ and $v$ that have coordinates $\begin{pmatrix} 3 \\ 2 \end{pmatrix}$ and $\begin{pmatrix} -2 \\ 6 \end{pmatrix}$ respectively? And in addition, why is that $6x+2y=22u$? Where is that coming from? Jul 7, 2018 at 9:14 • You have a parallelogram in the $(x,y)$ plane. It is a pain to integrate a function over this parallelogram. To make your life easier, consider a unit square $[0,1]^2$ in the $(u,v)$ plane. Via the linear transformation $$\begin{bmatrix} x\\ y\end{bmatrix} = \begin{bmatrix} 3 & -2\\ 2 & 6\end{bmatrix} \begin{bmatrix} u\\ v\end{bmatrix}$$ the $4$ vertices of the unit square in $(u,v)$ are mapped to the $4$ vertices of the parallelogram. Fortunately, the integrand in terms of $u$ and $v$ depends only on $u$. Thus, perform the integration over the unit square of the $(u,v)$ plane. Jul 7, 2018 at 12:46 • Note that the unit square in the $(u,v)$ plane is ($22$ times) smaller than the parallelogram in the $(x,y)$ plane. Thus, to compensate for the shrinking, multiply the integral over $u$ and $v$ by the determinant of the matrix (which is $22$). Jul 7, 2018 at 13:00
2022-08-10T08:54:48
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https://mathhelpforum.com/threads/family-of-three-children-conditional-probability.282503/
# Family of Three Children - Conditional Probability #### joshuaa From families with three children, a family is selected at random and found to have a boy. What is the probability that the boy has (a) an older brother and a younger sister; (b) an older brother; (c) a brother and a sister? Assume that in a three-child family all gender distributions have equal probabilities. #### Plato MHF Helper From families with three children, a family is selected at random and found to have a boy. What is the probability that the boy has (a) an older brother and a younger sister; (b) an older brother; (c) a brother and a sister? Assume that in a three-child family all gender distributions have equal probabilities. Here is the probability space: $\begin{array}{*{20}{c}}B&B&B\\B&B&G\\B&G&B\\B&G&G\\G&B&B\\G&B&G\\G&G&B\\G&G&G\end{array}$ You are given that the family has a boy. That creates a subspace. Now we just count. a) is $\dfrac{1}{7}$. Please explain that. #### joshuaa First, the book answer for (a) is 1/14 ------------------------- But let me try to tell you why you got 1/7 we have 8 outcomes in the sample space the probability of selecting a boy is 7/8 since only G G G does not have a boy B B G and B G B and G B B has two boys and a girl so the probability of getting two boys and a girl is 3/8 but if order is important only B B G will be the outcome of a boy has an older brother and a sister so the probability is 1/8 Therefore, the answer for (a) is (1/8) / (7/8) = 1/7 Is the answer of the book wrong? #### Plato MHF Helper First, the book answer for (a) is 1/14 we have 8 outcomes in the sample space the probability of selecting a boy is 7/8 since only G G G does not have a boy B B G and B G B and G B B has two boys and a girl so the probability of getting two boys and a girl is 3/8 but if order is important only B B G will be the outcome of a boy has an older brother and a sister so the probability is 1/8 Therefore, the answer for (a) is (1/8) / (7/8) = 1/7 Is the answer of the book wrong? If someone tell you that a family has three children and at least one is a boy, what is the probability of BBG? clearly order does matter. What is the probability of two boys and a girl? (order does not matter) Yes, I argue the answer is 1/7 as written. I would like to hear what that author says. #### romsek MHF Helper First, the book answer for (a) is 1/14 ------------------------- But let me try to tell you why you got 1/7 we have 8 outcomes in the sample space the probability of selecting a boy is 7/8 since only G G G does not have a boy B B G and B G B and G B B has two boys and a girl so the probability of getting two boys and a girl is 3/8 but if order is important only B B G will be the outcome of a boy has an older brother and a sister so the probability is 1/8 Therefore, the answer for (a) is (1/8) / (7/8) = 1/7 Is the answer of the book wrong? I think it's $\dfrac{1}{14}$ because getting GBB could be either the "boy" having an older brother as we want, or it could be that the "boy" is the older brother and thus has a younger brother. Both are equiprobable and they total to $\dfrac{1}{7}$ in probability so they are both $\dfrac{1}{14}$ in probability. #### Plato MHF Helper From families with three children, a family is selected at random and found to have a boy. What is the probability that the boy has (a) an older brother and a younger sister; (b) an older brother; (c) a brother and a sister? Assume that in a three-child family all gender distributions have equal probabilities. I think it's $\dfrac{1}{14}$ because getting GBB could be either the "boy" having an older brother as we want, or it could be that the "boy" is the older brother and thus has a younger brother. Both are equiprobable and they total to $\dfrac{1}{7}$ in probability so they are both $\dfrac{1}{14}$ in probability. @Romsek, it clearly says THAT boy has an older brother and younger sister. So the elementary event is BBG. We cannot count GBB because it violates the given. Last edited: #### romsek MHF Helper @Romsek, it clearly says THAT boy has an older brother and younger sister. So the elementary event is BBG. We cannot count GBB because it violates the given. No. It asks what is the probability that the boy has an older brother and younger sister. #### joshuaa GBB could be either the "boy" having an older brother as we want, or it could be that the "boy" is the older brother and thus has a younger brother. Both are equiprobable and they total to $\dfrac{1}{7}$ in probability so they are both $\dfrac{1}{14}$ in probability. I agree GBB are equiprobable, but they violate the event as Plato said Why not saying the Event BBG could be either the "boy" having an older brother as we want, or it could be that the "boy" is the older brother and thus has a younger brother. then 1/14 answer makes more sense Last edited: #### romsek MHF Helper I agree GBB are equiprobable, but they violate the event as Plato said Why not saying the Event BBG could be either the "boy" having an older brother as we want, or it could be that the "boy" is the older brother and thus has a younger brother. then 1/14 answer makes more sense That's what I meant. we know that BBG has probability $\dfrac 1 7$ BBG can be either our boy has an older brother, or that he is the older brother. Both with equal probability, and summing to $\dfrac 1 7$, i.e. $\dfrac{1}{14}$ So $P[\text{our boy has older brother and younger sister}] = \dfrac{1}{14}$ #### Plato MHF Helper From families with three children, a family is selected at random and found to have a boy. What is the probability that has (a) an older brother and a younger sister; (b) an older brother; (c) a brother and a sister? Assume that in a three-child family all gender distributions have equal probabilities. If I were tasked with editing this question, here is what I would submit. In a family of with three children what is the probability that a boy who has: A) an older brother and a younger sister? B) an older brother? C) a brother and a sister? The advantage in that wording is that any hint of a conditional is removed.
2020-01-26T15:39:17
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https://www.physicsforums.com/threads/differential-equation-proof.106673/
# Homework Help: Differential Equation Proof 1. Jan 13, 2006 I need some help. A pot of boiling water at 100C is removed from a stove at time t = o and left to cool in the kitchen. After 5 min, the water temperature has decreased to 80C, and another 5min later it has dropped to 65C. Assume Newtons law of cooling (dT/dt = k(M - T) ) applies, determain the constant temperature (M) of the kitchen. Okay, so the solution to the differential equation is - T= Ce^(kt) + M and we want to solve for M T(0) = C +M 100 = C + M Equation 1. T(5) = Ce^(k5) + M 80 = Ce^(k5) + M Equation 2. T(10) = Ce^(k10) + M 65 = Ce^(k10) + M Therefore.. the 2 equations are 80 = Ce^(k5) + M 65 = Ce^(k10) + M and if I subtract them, I get 15 = Ce^(k5) - Ce^(k10) And this is where im lost. I have no idea how to solve for M Thanks 2. Jan 13, 2006 ### saltydog So you have: $$\frac{dT}{dt}=k(M-T)$$ or: $$\frac{dT}{M-T}=kdt$$ Integrating: $$-ln(M-T)=kt+C$$ or: $$ln(M-T)=C-kt$$ (minus c, plus c same dif it's arbitrary) so that: $$M-T=e^{C-kt}$$ or: $$T(t)=M-ce^{-kt}$$ So you have 3 conditions, and 3 unknowns, kinda' messy but try and make those substitutions and then solve for something like M in terms of c, c in terms of k, then substitute that into one of them to solve for k, then the others. Last edited: Jan 13, 2006 3. Jan 13, 2006 ### HallsofIvy You are given the temperature at 3 times because you have 3 unknown numbers (C, k, and M) and need to solve 3 equations. You wrote the third equation but didn't number it or use it: C+ M= 100. You can use that to write the two equations 80= Ce5k+ M and 65= Ce10k+ M as 80= (100-M)e5k+ M and 65= (100-M)e10k+ M. Now, you want to eliminate k, not M, from those equations. Since e10k= (e5k)2 one way to do that is to write (80- M)2= (100-M)2e10k and 65- M= (100-M)e10k and divide one equation by the other. You get $$\frac{(80-M)^2}{65-M}= 100- M$$. That's a quadratic equation for M. 4. Jan 14, 2006 Salty dog, when i originally did solved the differential equation, i got the same answer as you, T= M - Ce^(-kt) However, my text book solution gives T= Ce^(kt) + M. Now im really confused. which one is right? Well... I did land at the same answer for M for both answers, so i assume they are both correct Last edited: Jan 14, 2006 5. Jan 14, 2006 ### Integral Staff Emeritus Since C, k, and M are arbitrary constants their algebraic sign is arbitrary until you have determined their value. 6. Jan 14, 2006 ### saltydog Ok guys. Didn't notice that. They are the same equation when solved for the constants: $$T=M-Ce^{-kt};\;M=20,\;C=-80,\;k=0.0575$$ $$T=M+Ce^{kt};\;M=20,\;C=80,\;k=-0.0575$$ 7. Jan 14, 2006 okay i have one more differential equation.... m(dv/dt) = mg - bv where m is the mass, g is the acceleration of gravity, and b>0 is a constant. if m = 100kg, g = 9.8m/s^2 and b=5 kg/sec, and v(0) = 10m/s, solve for v(t). what is the terminal velocity of the object. So firstly, i want to solve for this differential equation. But i cannot find a way to seperate it. (i mean, its not seperable?) m(dv/dt) = mg - bv (dv/dt) = g - bv/m then i pluged in the numbers and got (dv/dt) = 9.8 - 0.05v and thats about as far as i could get. could somebody please give me a little push thanks 8. Jan 14, 2006 ### qbert It is seperable, I'm not sure that's the way i would solve it but: (dv/dt) = g - (b/m)v => dv/(g - vb/m)= dt then a change of variables v -> vb/m turns this into a simple integral. 9. Jan 14, 2006 oh actually you are right... dv/(g - vb/m)= dt and you sub in the constants to get dv/(9.8-0.05v) = dt and solve (-1/0.05)ln|9.8 - 0.05v| = t + c and isolate for v to get v(t) = 196 - 20e^(-0.05t)(c) and given that v(0) = 10, you can solve for C to get v(t) = 196 - 186e^(-0.05t) Now.. the problem asks "What is the limiting (i.e, terminal) velocity of the object?" Would that just be the limit t-> infinity ? that would get 10m/s? 10. Jan 14, 2006 ### qbert the idea is right: v_term = limit (t->infinity) 196 - 186 exp(- t/20) = 196. given v(0) = v0 and m dv/dt = mg - bv we get: v(t) = gm/b + (v0 - gm/b)exp(-bt/m) and v_terminal = limit (t-> infinity) v(t). thus v_terminal = gm/b [ as lim t->infinity exp(-k*t) = 0 if k>0 ] we can check the answer, since at v_terminal there is no acceleration. that is dv/dt = 0 at v_terminal and the equation of motion is: m dv/dt = 0 = m g - b v_terminal => v_terminal = mg/b. i always like problems that have built in checks to see whether or not i get the right answer.
2018-04-23T17:42:10
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https://www.r-bloggers.com/2021/12/sample-and-population-variance-in-r/
Want to share your content on R-bloggers? click here if you have a blog, or here if you don't. Visit for the most up-to-date information on Data Science, employment, and tutorials finnstats. If you want to read the original article, go here Sample and Population Variance in R Sample and Population Variance in R, The variance is a metric for determining how dispersed data values are around the mean. Variance is the expectation of a random variable’s squared departure from its mean in probability theory and statistics, and it informally indicates how far a set of (random) values is spread out from its mean. How to Use the scale() Function in R » finnstats The formula for calculating a population’s variance is σ2 = Σ (xi – μ)2 / N where μ is the population mean, xi is the ith population element, N is the population size, and is basically Σ a fancy symbol for “sum.” To determine a sample’s variance, use the following formula: s2 = Σ (xi – xbar)2 / (n-1) where xbar represents the sample mean, xi represents the sample’s ith element, and n represents the sample size. ## Calculate Sample & Population Variance in R Assume we have the following R dataset and stored in data1. Cluster Analysis in R » Unsupervised Approach » finnstats Let’s create a data set values data1<- c(12,84, 5, 17, 18, 11, 13, 19, 69, 92,15,10,55) The var() function in R can be used to calculate sample variance. Let’s calculate the sample variance var(data1) 957.8974 The population variance can be calculated by multiplying the sample variance by (n-1)/n as follows. Now we can calculate the length of the data1 n <- length(data1) n 13 It’s ready to find population variance var(data1) * (n-1)/n 884.213 It’s important to remember that the population variance is always lower than the sample variance. Goodness of Fit Test- Jarque-Bera Test in R » finnstats In practice, we calculate sample variances for datasets because collecting data for a whole population is uncommon. Calculate the Sample Variance of Multiple Columns as an example Let’s say we have the following R data frame: Now we can create a data frame data2 <- data.frame(X=c(12, 35, 55, 48, 54, 12, 8, 10), Y=c(12, 24, 33, 77, 5, 46, 71, 106), Z=c(1, 2, 63, 8, 12, 77, 92, 102)) data2 X   Y   Z 1 12  12   1 2 35  24   2 3 55  33  63 4 48  77   8 5 54   5  12 6 12  46  77 7  8  71  92 8 10 106 102 To determine the sample variance of each column in the data frame, we can use the sapply() function: Yes, now based on sapply we can find each column’s sample variance. Regression Analysis Example-Ultimate Guide » finnstats sapply(data2, var) X         Y         Z 439.6429 1238.7857 1863.9821 We can also determine the sample standard deviation of each column using the following code, which is essentially the square root of the sample variance: To find each column’s sample standard deviation sapply(data2, sd) X        Y        Z 20.96766 35.19639 43.17386 When it comes to data analysis, Sapply is a highly handy function. Stringr in r 10 data manipulation Tips and Tricks » finnstats Don't forget to express your happiness by leaving a comment. Sample and Population Variance in R. If you are interested to learn more about data science, you can find more articles here finnstats. The post Sample and Population Variance in R appeared first on finnstats.
2022-01-19T19:06:55
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https://stats.stackexchange.com/questions/9825/changing-null-hypothesis-in-linear-regression/9827
# Changing null hypothesis in linear regression I have some data that is highly correlated. If I run a linear regression I get a regression line with a slope close to one (= 0.93). What I'd like to do is test if this slope is significantly different from 1.0. My expectation is that it is not. In other words, I'd like to change the null hypothesis of the linear regression from a slope of zero to a slope of one. Is this a sensible approach? I'd also really appreciate it you could include some R code in your answer so I could implement this method (or a better one you suggest!). Thanks. ## 5 Answers set.seed(20); y = rnorm(20); x = y + rnorm(20, 0, 0.2) # generate correlated data summary(lm(y ~ x)) # original model summary(lm(y ~ x, offset= 1.00*x)) # testing against slope=1 summary(lm(y-x ~ x)) # testing against slope=1 Outputs: Estimate Std. Error t value Pr(>|t|) (Intercept) 0.01532 0.04728 0.324 0.75 x 0.91424 0.04128 22.148 1.64e-14 *** Estimate Std. Error t value Pr(>|t|) (Intercept) 0.01532 0.04728 0.324 0.7497 x -0.08576 0.04128 -2.078 0.0523 . Estimate Std. Error t value Pr(>|t|) (Intercept) 0.01532 0.04728 0.324 0.7497 x -0.08576 0.04128 -2.078 0.0523 . • Thank you! I just couldn't figure out how to change the lm command. Apr 21, 2011 at 14:35 • Then is it exactly the same "lm(y-x ~ x)" than "lm(y ~ x, offset= 1.00*x)" (or without that 1.00) ? Wouldn't that substraction be problem with the assumptions for least squares or with collinearity? I want to use it for a logistic regression with random effects glmer(....). It would be great to have a simple but correct method to get the p-values. – skan May 19, 2017 at 15:43 • Here stats.stackexchange.com/questions/111559/… Matifou says this method is worse than using Wald the test. – skan May 19, 2017 at 15:45 Your hypothesis can be expressed as $R\beta=r$ where $\beta$ is your regression coefficients and $R$ is restriction matrix with $r$ the restrictions. If our model is $$y=\beta_0+\beta_1x+u$$ then for hypothesis $\beta_1=0$, $R=[0,1]$ and $r=1$. For these type of hypotheses you can use linearHypothesis function from package car: set.seed(20); y = rnorm(20); x = y + rnorm(20, 0, 0.2) # generate correlated data mod <- lm(y ~ x)) # original model > linearHypothesis(mod,matrix(c(0,1),nrow=1),rhs=c(1)) Linear hypothesis test Hypothesis: x = 1 Model 1: restricted model Model 2: y ~ x Res.Df RSS Df Sum of Sq F Pr(>F) 1 19 0.96022 2 18 0.77450 1 0.18572 4.3162 0.05234 . --- Signif. codes: 0 ‘***’ 0.001 ‘**’ 0.01 ‘*’ 0.05 ‘.’ 0.1 ‘ ’ 1 • Can this be used for a one-sided test? – user66430 Feb 7, 2017 at 15:17 It seems you're still trying to reject a null hypothesis. There are loads of problems with that, not the least of which is that it's possible that you don't have enough power to see that you're different from 1. It sounds like you don't care that the slope is 0.07 different from 1. But what if you can't really tell? What if you're actually estimating a slope that varies wildly and may actually be quite far from 1 with something like a confidence interval of ±0.4. Your best tactic here is not changing the null hypothesis but actually speaking reasonably about an interval estimate. If you apply the command confint() to your model you can get a 95% confidence interval around your slope. Then you can use this to discuss the slope you did get. If 1 is within the confidence interval you can state that it is within the range of values you believe likely to contain the true value. But more importantly you can also state what that range of values is. The point of testing is that you want to reject your null hypothesis, not confirm it. The fact that there is no significant difference, is in no way a proof of the absence of a significant difference. For that, you'll have to define what effect size you deem reasonable to reject the null. Testing whether your slope is significantly different from 1 is not that difficult, you just test whether the difference $slope - 1$ differs significantly from zero. By hand this would be something like : set.seed(20); y = rnorm(20); x = y + rnorm(20, 0, 0.2) model <- lm(y~x) coefx <- coef(summary(model))[2,1] seslope <- coef(summary(model))[2,2] DF <- model$df.residual # normal test p <- (1 - pt(coefx/seslope,DF) )*2 # test whether different from 1 p2 <- (1 - pt(abs(coefx-1)/seslope,DF) )*2 Now you should be aware of the fact that the effect size for which a difference becomes significant, is > qt(0.975,DF)*seslope [1] 0.08672358 provided that we have a decent estimator of the standard error on the slope. Hence, if you decide that a significant difference should only be detected from 0.1, you can calculate the necessary DF as follows : optimize( function(x)abs(qt(0.975,x)*seslope - 0.1), interval=c(5,500) )$minimum [1] 6.2593 Mind you, this is pretty dependent on the estimate of the seslope. To get a better estimate on seslope, you could do a resampling of your data. A naive way would be : n <- length(y) seslope2 <- mean( replicate(n,{ id <- sample(seq.int(n),1) model <- lm(y[-id]~x[-id]) coef(summary(model))[2,2] }) ) putting seslope2 in the optimization function, returns : \$minimum [1] 6.954609 All this will tell you that your dataset will return a significant result faster than you deem necessary, and that you only need 7 degrees of freedom (in this case 9 observations) if you want to be sure that non-significant means what you want it means. You can simply not make probability or likelihood statements about the parameter using a confidence interval, this is a Bayesian paradigm. What John is saying is confusing because it there is an equivalence between CIs and Pvalues, so at a 5%, saying that your CI includes 1 is equivalent to saying that Pval>0.05. linearHypothesis allows you to test restrictions different from the standard beta=0
2022-05-25T14:15:55
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http://math.stackexchange.com/questions/53950/how-many-hands-to-reduce-the-probability-of-no-royal-flush-in-n-hands-to-less-th
# How many hands to reduce the probability of no royal flush in n hands to less than 1/e? I get close but can't figure out how they get the answer to this question -- Introduction to Probability, Grinstead, Snell, Chapter 5.1, exercise 17: The probability of a royal flush in a poker hand is p = 1/649,740. How large must n be to render the probability of having no royal flush in n hands smaller than 1/e? My answer is: since this is for one success, use a geometric distribution instead of a negative binomial. • success p(flush) = 1/649740 = .0000016 • failure p(no flush) = 1 - 1/649740 = .9999984 • 1/e = 1 / 2.71828 = .36788 For geometric distribution, I set them equal to find the minimum n: • q^n = 1/e • (.9999984)^n = .36788 • log _.9999984 (.36788) = n (base is .9999984, the failure probability) • log .36788 / log .9999984 = n (found this technique for using base 10 for both) • n = 624998.55 Therefore, my answer is 624999 (just over the equal n by rounding up). However, the answer in the text book says 649741. Using my above process, I get that too -- if I change 1/e from .36788 to .3536. Is my 1/e incorrect? Or have I used the wrong distribution and formula? Thanks in advance for any help. - I hope that when you use your calculator, you do not truncate numbers and feed the truncated version back in. Use the memory feature! It will save on time and on keying errors. Also, calculators do their calculation internally to higher precision than they reveal externally. If you key something back in, even to "full precision", you are losing accuracy. And you do not need logs to weird bases. Every additonal step in a calculation is a threat to accuracy. – André Nicolas Jul 27 '11 at 1:17 thanks for the tip -- I was doing it the wrong way as you note, I'll follow these steps next time. I need a new calculator, the one I have does not have an e button. – d l Jul 27 '11 at 23:19 So you don't have $\ln x$, $e^x$ on your calculator? Your computer, whatever operating system it runs, almost certainly has a scientific calculator utility. If it doesn't, a free one can be downloaded. – André Nicolas Jul 28 '11 at 2:02 Oh no -- I'm getting more embarrassed; I have the calculator in Windows 7, Scientific View, and it has "ln", but no "e" .... until I hit the "Inv" button, the function buttons change and "e" is there. Sorry for asking a dumb question, and many thanks for your persistence, otherwise I would not have explored further and found this. – d l Jul 28 '11 at 18:27 Nice design, I guess, but too clever. On physical calculators, $e^x$ is permanently written under the key. – André Nicolas Jul 28 '11 at 19:12 This is just an answer due to rounding. You correctly solve the equation $$q^n = 1/e$$ for $n$ to get: \begin{align} n = \frac{\log(1/e)}{\log(q)} = \frac{-1}{\log(q)} \approx 649739.5715, \end{align} according to my calculator. Notice that there is no need to explicitly solve for $e$ since $\log(1/e) = -1$. - thanks -- I did not know that about e, I think that's where my problem was. – d l Jul 27 '11 at 18:42 You're quite welcome. That just comes from $\log_b(b) = 1$. Indeed, when a mathematician writes $\log$ (s)he means the natural log $\ln = \log_e$. Since all logs are the same (up to constants), we just stick with the notation $\log$. – JavaMan Jul 27 '11 at 18:45 It's a rounding error. You can get closer results as follows: \begin{align} q^n &= e^{-1} \\ n \ln q &= -1 \\ n &= 1 / (-\ln q) \end{align} Note that $-\ln(1 - x)$ is approximately $x$ for small x. Just plugging in - 1/649740 gives a result of 649740. The actual series is $x + x^2 / 2 + x^3/3 ...$, which means that the $-\ln q > 1/649740$, but only very slightly, which means that $-1/\ln q < 649740$, so 649740 is sufficient, rather than the 649741 that the book claims. - thanks -- I get the same answer with your solution if I hit "ln" on my calculator, not "log". I was using "log" before, and did not know e inverse = -1. – d l Jul 27 '11 at 18:45 When $p$ is small, the probability of at least one success in $n$ trials is approximately $1-1/\mathrm e$ when $n = 1/p$. Your $p = 1/649740$ should be plenty small enough for this approximation to be very good. To be exact, the probability of no success in $n$ independent trials, with probability $p$ of success in each, is $(1-p)^n = \exp(n \log(1-p))$. When $n$ is large and/or $p$ is small, this is quite well approximated by $\exp(-np)$. Equivalently, writing $k = np$ for the expected number of successes in the $n$ trials, the probability of no success is approximately $\exp(-k)$ regardless of $n$ (as long as $n$ isn't very small). -
2016-05-03T19:06:31
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https://math.stackexchange.com/questions/1283195/can-we-determinine-the-convergence-of-int-0-infty-fracx2n-1x2-1
# Can we determinine the convergence of $\int_0^\infty \frac{x^{2n - 1}}{(x^2 + 1)^{n + 3}}\,dx$ without evaluating it? Can we determine convergence without evaluating this improper integral? $$\int_0^\infty {\frac{x^{2n - 1}}{{\left( x^2 + 1 \right)}^{n + 3}}\,dx}\quad\quad n\geq 1\;,\; n\in\mathbb{Z}$$ When trying to bound the integrand, relating a known function, the integral does not converge. But when I use the software package Maple I see that the integral converges to: $$\int_0^\infty {\frac{x^{2n - 1}}{{\left( x^2 + 1 \right)}^{n + 3}}\,dx}=\frac{1}{{\left( {n + 1} \right)\left( {n + 2} \right)n}}$$ Evaluating it is not difficult, but I want to know if you can determine convergence (delimiting the integrand) without evaluating such an improper integral. • Comparison test is the first thing to pop into my mind. – Vim May 15 '15 at 10:04 • Show that $0\leq \frac{x^{2n-1}}{\left(x^2+1\right)^{n+3}}\leq \frac{1}{x^2+1}$ for $x>0$ and all $n\geq 1$ – Lozenges May 15 '15 at 11:44 • The integrand is asymptotically $x^{2n-1-2(n+3)}=x^{-7}$, so convergence is very fast. – Yves Daoust May 15 '15 at 14:10 • @Lozenges I find your suggestion very useful, but this inequality is difficult to prove. Anyway thanks. – mathsalomon May 15 '15 at 19:39 • quite easy to prove by induction on $n$ that $0\leq x^{2n-1}\leq \left(x^2+1\right)^{n+2}$ – Lozenges May 16 '15 at 8:57 By equivalents: $$\frac{x^{2\;n - 1}}{(x^2 + 1)^{n + 3}}\sim_{\infty}\frac{x^{2\;n - 1}}{x^{2(n + 3)}}=\frac1{x^7}$$ and the latter is convergent since the exponent of $x$ in the denominator is $>1$. • And close to $x=0$, the integrand is $x^{2 n} \left(\frac{1}{x}-(n+3) x+O\left(x^3\right)\right)$ – Claude Leibovici May 15 '15 at 8:59 • The hypothesis is $n\ge1$, so there is no problem of convergence. – Bernard May 15 '15 at 9:06 • I noticed but I prefered to control ! Cheers :-) – Claude Leibovici May 15 '15 at 9:09 • @Bernard Thank you, just do not know how to justify to $x \in [0,1]$. Thanks to Ian comment now know how to justify it. – mathsalomon May 15 '15 at 19:45 • For $x\in [0,1]$, there's nothing to justify since the integrand is continuous on $[0,1]$, so it is an ordinary Riemann integral. – Bernard May 15 '15 at 20:02 The numerator is of degree $2n - 1$. The denominator looks like a polynomial of degree $2n + 6$, and is bounded away from $0$. As a general rule, integrals of the form $$\int_0^\infty \frac{1}{(1 + x)^n}dx$$ converge for $n > 1$ and diverge for $n \leq 1$. Your integral looks like this integral, but with $n = 7$. So it converges. You can make this argument more formal through limit comparison, if that's something you're familiar with. The naive comparison would simply drop the $1$ in the denominator, but the result is not bounded, whereas your integrand is bounded. To fix this, keep the $1$ on, say, $[0,1]$ and then drop it on $[1,\infty)$. The effect is that on $[0,1]$ the integrand is bounded while on $[1,\infty)$ the integrand is bounded, and in particular bounded by $Cx^{-7}$ for some $C>0$. • Thanks, that's just what he had considered to do, because of the divergence of $x^{-7}$ in $[0,1]$. – mathsalomon May 15 '15 at 19:31 Setting $x=\tan(t)$, we obtain $$\dfrac{x^{2n-1}}{\left(x^2+1\right)^{n+3}} = \dfrac{\tan^{2n-1}(t)}{\sec^{2n+6}(t)} = \sin^{2n-1}(t)\cos^7(t)$$ Hence, we have $$\int_0^{\infty}\dfrac{x^{2n-1}}{\left(x^2+1\right)^{n+3}} dx = \int_0^{\pi/2} \sin^{2n-1}(t)\cos^5(t)dt$$ which clearly exists, since the integrand is continuous and bounded on $[0,\pi/2]$. Also, what better way to prove convergence than evaluating the integral. The integral can be easily evaluated using the identity here and the definition of the $\beta$ function. We have $$\beta(x,y) = 2 \int_0^{\pi/2}\sin^{2x-1}(t)\cos^{2y-1}(t)dt$$ which gives us that $$\int_0^{\pi/2} \sin^{2n-1}(t)\cos^5(t)dt = \dfrac{\beta(n,4)}2 = \dfrac{\Gamma(n)\Gamma(3)}{2\Gamma(n+3)} = \dfrac1{(n+2)(n+1)n}$$ • Did you manage the $dx$? (I think you should have another $\sec^2(t)$.) – Ian May 15 '15 at 14:23 • @Ian Thanks. corrected. – Leg May 15 '15 at 14:28
2019-07-18T19:08:58
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http://zonasantos.com/bohai-sea-pflcbgm/sigma-notation-examples-e0b873
1. Let x 1, x 2, x 3, …x n denote a set of n numbers. If f(i) represents some expression (function) involving i, then has the following meaning : . Active 6 years, 10 months ago. (By the way: The summation formula can be proved using induction.). Click HERE to return to the list of problems. It’s just a “convenience” — yeah, right. In the content of Using Sigma Notation to represent Finite Geometric Series, we used sigma notation to represent finite series. (2 answers) Closed 6 years ago. We’ll start out with two integers, $$n$$ and $$m$$, with $$n < m$$ and a list of numbers denoted as follows, Properties . Really clear math lessons (pre-algebra, algebra, precalculus), cool math games, online graphing calculators, geometry art, fractals, polyhedra, parents and teachers areas too. Watch Queue Queue. Set-Builder Notation. The summation notation is a way to quickly write the sum of a series of functions. Download fifa 13 soundtrack Messages. Section 7-8 : Summation Notation. Instead of using the f(x) notation, however, a sequence is listed using the a n notation. An infinity symbol ∞ is placed above the Σ to indicate that a series is infinite. Worked examples summation notation. T HIS —Σ—is the Greek letter sigma. Unsure of sigma notation. The concept of sigma notation means to sum up all terms and uses three parts to form math statements, like ∑ i a i.The Greek letter ∑ is the summation operator and means the sum of all, i is called the index number, and a i refers to a series of terms to be added together. In this section we need to do a brief review of summation notation or sigma notation. In this case we'd think of the general term as Use summation notation to write the series. This video is unavailable. The break point is usually obvious from standard rules for algebraic expressions, or other aspects of the notation, In order to master the techniques explained here it is vital that you undertake plenty of practice exercises so that they become second nature. The summation operator governs everything to its right. sigma notation, also known as summation notation. {x : x > 0} means "the set of all x such that x is greater than 0". The definition implies that it also includes the empty subset and that it is closed under countable intersections.. Cross your fingers and hope that your teacher decides not […] Description. A shorthand used to write sets, often sets with an infinite number of elements. The "X i" indicates that X is the variable to be summed as i goes from 1 to 4. For example: This means that we are to repeatedly add ka k. The first time we write it, we put k = 1. The sum of consecutive numbers. Remainder classes modulo m. An arithmetic series. *Please select more than one item to compare In this unit we look at ways of using sigma notation, and establish some useful rules. It is used like this: Sigma is fun to use, and can do many clever things. This mathematical notation is used to compactly write down the equations in which summing all terms is required. The variable Search results for download at Sigma-Aldrich. The following diagram shows the Sigma Notation. Example 1.1 . Demonstrates how to find the value of a term from a rule, how to expand a series, how to convert a series to sigma notation, and how to evaluate a recursive sequence. Sigma notation. Go To Problems & Solutions Return To Top Of Page . 7.1 - Sequences and Summation Notation. Alternatively, we could decide we wanted to write the series starting at n = 0. Thinking of the summation formula this way can be a useful way of memorizing the formula. Learn more at Sigma Notation.. You might also like to read the more advanced topic Partial Sums.. All Functions The index of summation , here the letter i, is a dummy variable whose value will change as the addends of the sum change. SOLUTIONS TO THE ALGEBRA OF SUMMATION NOTATION SOLUTION 1 : = (5+1) + (5+2) + (5+4) + (5+8) = 6 + 7 + 9 + 13 = 35 . Stress's. SUMMATION (SIGMA) NOTATION 621 Getting back to this particular proof, the statement P1 would be that 1 X i3 = i=1 11 (1 + 1)2 , 4 2 2 which is clearly true because it is equivalent to 13 = 1 (2) 4 , i.e., 1 = 1, which is true (obviously). Proof . Summation notation is used to define the definite integral of a continuous function of one variable on a closed interval. Often mathematical formulae require the addition of many variables Summation or sigma notation is a convenient and simple form of shorthand used to give a concise expression for a sum of the values of a variable. Properties of Sigma Notation - Cool Math has free online cool math lessons, cool math games and fun math activities. The "i = 1" at the bottom indicates that the summation is to start with X 1 and the 4 at the top indicates that the summation will end with X 4. 2. 1. Notation . Dismantled. Sigma notation examples with answers. Write out these sums: Solution. By the way, you don’t need sigma notation for the math that follows. EOS . Wettest. Scroll down the page for more examples and solutions using the Sigma Notation. A typical sum written in sigma notation looks like this: 4 k 0 (k2 3) The symbol “Σ” is the Greek capital letter sigma, which stands for “sum”. Summation Notation And Formulas . The letter sigma is a signal that summation notation is being used. It may also be any other non-negative integer, like 0 or 3. Series : Sigma Notation : ExamSolutions : A-Level Maths In this tutorial you are shown the meaning behind sigma notation for the sum of a sequence called a series. Provides worked examples of typical introductory exercises involving sequences and series. Sigma Notation. Summation notation uses the sigma Σ symbol to represent sums with multiple terms. Shows how factorials and powers of –1 can come into play. There are infinite sequences whose domain is the set of all positive integers, and there are finite sequences whose domain is the set of the first n positive integers. Search results for msds at Sigma-Aldrich. [duplicate] Ask Question Asked 6 years, 10 months ago. Summation notation works according to the following rules. You can use sigma notation to write out the right-rectangle sum for a function. Beautiful lyrics download Download gta vice city 5 game free. The lower limit of the sum is often 1. $\endgroup$ – nbro Dec 19 '16 at 15:33 Let's first briefly define summation notation. explaining using examples how to overcome or try to overcome the difficulties in interpreting this notations. Compare Products: Select up to 4 products. Psychologists Sigma notation exercises. Worked examples: summation notation … Sepulchral. A sequence is a function whose domain is the natural numbers. Sigma notation uses a variable that counts upward to change the terms in the list. Riemann sums, summation notation, and definite integral notation Summation notation We can describe sums with multiple terms using the sigma operator, Σ. Sigma notation for sums topics in precalculus. The induction step (2) has a simple, yet sophisticated little proof. Summation notation. x+0=4 Simplify. I don't understand the sigma notation and for loop stack overflow. Hippies. The Sigma symbol, , is a capital letter in the Greek alphabet.It corresponds to “S” in our alphabet, and is used in mathematics to describe “summation”, the addition or sum of a bunch of terms (think of the starting sound of the word “sum”: Sssigma = Sssum). Moderately Seria facil luis fonsi download. Watch Queue Queue In mathematics, summation is the addition of a sequence of any kind of numbers, called addends or summands; the result is their sum or total. The sum of the first n terms of a series is called "the n-th partial sum", and is often denoted as "S n ". But with sigma notation (sigma is the 18th letter of the Greek alphabet), the sum is much more condensed and efficient, and you’ve got to admit it looks pretty cool: This notation just tells you to plug 1 in for the i in 5i, then plug 2 into the i in 5i, then 3, then 4, and so on all … The Greek letter capital sigma (Σ) indicates summation. SOLUTION 2 : (The above step is nothing more than changing the order and grouping of the original summation.) up to a natural break point in the expression. Therefore, Sigma (Summation) Notation. Sigma notation examples. x 1 is the first number in the set. For example, say you’ve got f (x) = x2 + 1. Return To Contents Go To Problems & Solutions . Snowmobiles. Summation notation solutions. 5(0.3) 5 + 5(0.3) 6 + 5(0.3) 7 + .... Then we would write the series as. Viewed 4k times 1 $\begingroup$ This question already has answers here: Induction proof that $\sum_{j=n}^{2n-1} (2j + 1) = 3n^2$ - what happened? Sigma notation mc-TY-sigma-2009-1 Sigma notation is a method used to write out a long sum in a concise way. In mathematical analysis and in probability theory, a σ-algebra (also σ-field) on a set X is a collection Σ of subsets of X that includes X itself, is closed under complement, and is closed under countable unions.. $\begingroup$ Not at the moment, but I would cheerfully read an article talking about the topic, i.e. SIGMA NOTATION FOR SUMS. Compare Products: Select up to 4 products. We use it to indicate a sum. *Please select more than one item to compare The pair (X, Σ) is called a measurable space or Borel space. That is indicated by the lower index of the letter 5(0.3) 5 + 5(0.3) 6 + 5(0.3) 7 + .... 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2021-12-09T14:43:20
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http://partshelf.com/m96py/generate-positive-definite-matrix-r-747b34
## generate positive definite matrix r As is always the case for the generation of random objects, you need to be careful about the distribution from which you draw them. Positive Definite Matrix. More specifically, we will learn how to determine if a matrix is positive definite or not. ($$\lambda_1,\ldots,\lambda_p$$) for the covariance matrix A symmetric matrix is defined to be positive definite if the real parts of all eigenvalues are positive. matrix<-.pdMat. Note that all José Pinheiro and Douglas Bates [email protected]. sqrt(Sii), i=1,..,n and The remaining methods, denoted as “onion”, “c-vine”, and “unifcorrmat” eigenvalues: Vector of eigenvalues desired in output. Following are papers in the field of stochastic precipitation where such matrices are used. If any of the eigenvalues in absolute value is less than the given tolerance, that eigenvalue is replaced with zero. Show Hide all comments. evaluated on a data.frame to resolve the names it defines. Symmetrisch positiv-definite Matrizen Satz 3.1. Satz 3.2. positive definite matrix/covariance matrix. This function is a constructor for the pdNatural class, representing a general positive-definite matrix, using a natural parametrization . matrix in natural parametrization, also inheriting from class The simplest to produce is a square matrix size(n,n) that has the two positive eigenvalues 1 and n+1. Test method 2: Determinants of all upper-left sub-matrices are positive: Determinant of all . ($$\boldsymbol{Q}=(\boldsymbol{\alpha}_1,\ldots,\boldsymbol{\alpha}_p)$$) uninitialized object. In our experience, lambdaLow$$=1$$ and ratioLambda$$=10$$ Lower bound on the eigenvalues of cluster covariance matrices. Your method will produce a matrix R that looks "like" a correlation matrix, but beware - it is an impostor! For a positive definite matrix, the eigenvalues should be positive. If value is an initialized pdMat object, as the dimension increases. The covariance matrix $$\boldsymbol{\Sigma}$$ is then $\begingroup$ I encounter the problem of not positive definite matrices Your second matrix (following these words) appears negatively definite. coef.pdMat, Hello I am trying to determine wether a given matrix is symmetric and positive matrix. assigned later, generally using the coef or matrix replacement May 19, 2013 at 9:47 pm: On Sun, May 19, 2013 at 5:04 PM, Gabor Grothendieck wrote: On Sun, May 19, 2013 at 4:57 PM, Gabor Grothendieck wrote: On Sun, May 19, 2013 at 10:33 AM, mary wrote: user-specified eigenvalues when covMethod = "eigen". Yes, that's possible. This definition makes some properties of positive definite matrices much easier to prove. I want to generate a positive definite matrix such that all the correlations have tighter than trivial bounds. NOT be used for optimization. p. 162. as.matrix.pdMat, Today, we are continuing to study the Positive Definite Matrix a little bit more in-depth. The Cartan matrix of a simple Lie algebra is the matrix whose elements are the scalar products = (,) (,) (sometimes called the Cartan integers) where r i are the simple roots of the algebra. May 19, 2013 at 2:31 pm: Hi, I have a question for my simulation problem: I would like to generate a positive (or semi def positive) covariance matrix, non singular, in wich the spectral decomposition returns me the same values for all dimensions but differs only in eigenvectors. B=A.^(1/2) %scale down to range 0-4. Sign in to comment. Defaults to Defaults to NULL. The matrix has real valued elements. I know A'A will give a symmetric positive definite matrix. In other words, for every positive number R and increment h, the k-element vector {R, R-h, R-2h, ..., R-(k-1)h} generates a valid covariance matrix provided that R-(k-1)h > 0, which is equivalent to h ≤ R/(k-1). 262 POSITIVE SEMIDEFINITE AND POSITIVE DEFINITE MATRICES Proof. [R] how to randomly generate a n by n positive definite matrix in R ? Method to generate positive definite matrices/covariance matrices. share | cite | … S(i,j)/sqrt(S(i,i)S(j,j)), i not equal to j denote the associated Let R be a symmetric indefinite matrix, that is, a matrix with both positive and negative eigenvalues. as an uninitialized pdSymm object (with just some of its 4. I.e. Value Sigma the covariance matrix A symmetric square root of Sigma shift how much the eigenvalues were shifted. Quellcode-Beispiel (Python): from scipy import random, linalg matrixSize = 10 A = random.rand(matrixSize,matrixSize) B = numpy.dot(A,A.transpose()) print 'random positive semi-define matrix for today is', B But its still better to produce a positive-definite covariance matrix in a principled way from some model. Break the matrix in to several sub matrices, by progressively taking . Joe, H. (2006) a matrix of class dpoMatrix, the computed positive-definite matrix. (not only for your responses in this email thread but in helping create R generally and many of these functions in particular.) thanks! pdMat. formula, or a vector of character strings, object is returned set.seed(1) n <- 10 ## Dimension of matrix m <- 1000 ## Number of samples ## Create sparse, symmetric PSD matrix S A <- rsparsematrix(n, n, 0.15, rand.x = stats::rnorm) Strue <- A %*% t(A) + 0.05 * diag(rep(1, n)) ## Force matrix to be strictly positive definite. Example-Prove if A and B are positive definite then so is A + B.) Finally, if value is a numeric If It is known that a positive definite matrix has a Unique Positive Definite square root. General Positive-Definite Matrix Description. interval [lambdaLow, lambdaLow$$*$$ratioLambda]. [R] Generate positive definite matrix with constraints; Gabor Grothendieck. Theorem C.6 The real symmetric matrix V is positive definite if and only if its eigenvalues Method to generate positive definite matrices/covariance matrices. Often such matrices are intended to estimate a positive definite (pd) matrix, as can be seen in a wide variety of psychometric applications including correlation matrices estimated from pairwise or binary information (e.g., Wothke, 1993). parametrization . alphad=1 for uniform. nnode: Number of nodes in the matrix. The code enclosed has created such a function that will create a positive definite matrix of any size n x n. user-specified eigenvalues when covMethod = "eigen". Eine reelle symmetrische quadratische Matrix = (,), = ist genau dann positiv definit, wenn das Gaußsche Eliminationsverfahren bei Diagonalstrategie, das heißt ohne Zeilenvertauschungen, mit n positiven Pivotelementen durchgeführt werden kann. I could generate the matrices using an uniform distribution (as far as I could see, this is the standard method) and then force it to be positive-definite using this. Because There are MANY issues here. as.data.frame.sparsebnData: Convert a sparsebnData object back to a data.frame as.edgeList: as.edgeList as.sparse: as.sparse coerce_discrete: Recode discrete data count.interventions: Count the number of rows under intervention count.levels: Count the number of levels per variable degrees: Degree distribution of a graph Dimension of the matrix to be generated. Generating Random Correlation Matrices Based on Partial Correlations. ACM Transactions on Modeling and Computer Simulation (TOMACS), The ratio of the upper bound of the eigenvalues to the lower bound of the If the matrix associated with object is of dimension n, it is represented by n*(n+1)/2 unrestricted parameters, using the matrix-logarithm parametrization described in Pinheiro and Bates (1996). correlation matrix ($$\boldsymbol{R}$$) via the method mentioned and proposed in Joe (2006), See lambdaLow. dimensions and the row/column names of the underlying matrix. Defaults to the A convenient choice is h = R / k. This is a useful fact because it enables you to construct arbitrarily large Toeplitz matrices from a decreasing sequence. A positive definite matrix will have all positive pivots. Eigenvalues of a positive definite real symmetric matrix are all positive. Next message: [R] independence of censoring in survival analyses Messages sorted by: unrestricted (meaning that not all unrestricted vectors would give $$diag(\sigma_1,\ldots,\sigma_p)*\boldsymbol{R}*diag(\sigma_1,\ldots,\sigma_p)$$. Next message: [R] independence of censoring in survival analyses Messages sorted by: NULL, no attempt is made to obtain information on This function computes the nearest positive definite of a real symmetric matrix. r(i,j) = If eigenvalue = NULL and covMethod = "eigen", then eigenvalues will be automatically generated. Range for variances of a covariance matrix (see details). +), a vector of character strings, or a numeric matrix and unreplicated elements. following: a pdMat object, a positive-definite It is used to Therefore, this parametrization should it has some negative eigenvalues (and no zero eigenvalues). positive-definite matrices). The value is in binary and indication is available on the number of binary places to move over. Here denotes the transpose of . eigenvalues of cluster covariance matrices. When value is In linear algebra, a symmetric × real matrix is said to be positive-definite if the scalar is strictly positive for every non-zero column vector of real numbers. The default range is $$[1, 10]$$ which can generate reasonable Solution method B finds the nearest (to the original matrix) positive definite matrix having the specified minimum eigenvalue, in the sense of minimum frobenius norm of the difference of the positive definite matrix D and the original matrix C, which is based on the sums of squared differences of all elements of D - C, to include the off-diagonal elements. Also, it is the only symmetric matrix. Kurowicka and Cooke, 2006. parent frame from which the function was called. attributes and its class defined) and needs to have its coefficients length equal to the dimension of the underlying positive-definite Generate a random positive definite matrix. Method to generate positive definite matrices/covariance matrices. [R] Generate positive definite matrix with constraints; Mary. [R] Generate positive definite matrix with constraints; Gabor Grothendieck. The R function eigen is used to compute the eigenvalues. Your method will produce a matrix R that looks "like" a correlation matrix, but beware - it is an impostor! If the matrix associated with object is of dimension n, it is represented by n*(n+1)/2 parameters. lambdaLow should be positive. eta=1 for uniform. A matrix is positive definite fxTAx > Ofor all vectors x 0. A=16*gallery('lehmer',100) %matrix of size 100*100 in range 0-16. If truly positive definite matrices are needed, instead of having a floor of 0, the negative eigenvalues can be converted to a small positive number. If eigenvalue = NULL and covMethod = "eigen", then eigenvalues will be automatically generated. If the matrix associated with object is of dimension n, it is represented by n*(n+1)/2 parameters. parameter for unifcorrmat method to generate random correlation matrix In linear algebra, a symmetric × real matrix is said to be positive-definite if the scalar is strictly positive for every non-zero column vector of real numbers. obtain the levels for factors, which affect the No real data (having no missings) can ever correspond to such a covariance matrix. – vak Jun 24 '09 at 12:27 | show 1 more comment. Sign in to answer this question. pdClasses, In that case, if S in the above decomposition is positive definite, then A is said to be a Cartan matrix. An integer in R consists of the whole number that can be positive or negative whereas a floating-point number includes real numbers. See help("make.positive.definite") from package corpcor. RDocumentation Transposition of PTVP shows that this matrix is symmetric.Furthermore, if a aTPTVPa = bTVb, (C.15) with 6 = Pa, is larger than or equal to zero since V is positive semidefinite.This completes the proof. There is more structure to a correlation matrix than that meets the eye! Unfortunately, with pairwise deletion of missing data or if using tetrachoric or polychoric correlations, not all correlation matrices are positive definite. The matrix symmetric positive definite matrix A can be written as , A = Q'DQ , where Q is a random matrix and D is a diagonal matrix with positive diagonal elements. an optional vector of character strings specifying the [R] Generate positive definite matrix with constraints; Mary. uses columns of a randomly generated orthogonal matrix Choices are “eigen”, “onion”, “c-vine”, or “unifcorrmat”; see details below. That... could work. Generate a random positive definite matrix Usage. Afterwards, the matrix is recomposed via the old eigenvectors and new eigenvalues, and then scaled so that the diagonals are all 1′s. num.ortho: Number of random Householder reflections to compose. Wiley, 2006. Because the diagonal is 1 and the matrix is symmetric. The eigenvalues are randomly generated from the normF: the Frobenius norm (norm(x-X, "F")) of the difference between the original and the resulting matrix. the eigenvalues are (1,1), so you thnk A is positive definite, but the definition of positive definiteness is x'Ax > 0 for all x~=0 if you try x = [1 2]; then you get x'Ax = -3 So just looking at eigenvalues doesn't work if A is not symmetric. matrix, a one-sided linear formula (with variables separated by If the argument covMethod="eigen", eigenvalues are generated for cluster covariance matrices. eigenvalue. Behavior of the NORTA method for correlated random vector generation upper-left sub-matrices must be positive. NULL. It is mostly used for deriving factors appearing in the formulas. In such cases one has to deal with the issue of making a correlation matrix positive definite. functions. However, I found that *Lehmer* matrix is a positive definite matrix that when you raise each element to a nonnegative power, you get a positive semi-definite matrix. I) dIiC fifl/-, Our final definition of positive definite is that a matrix A is positive definite if and only if it can be written as A = RTR, where R is a ma trix, possibly rectangular, with independent columns. numeric. of a positive definite matrix. methods to generate random covariance matrices. Ravi Varadhan rvaradhan at jhmi.edu Thu Feb 7 20:02:30 CET 2008. Choices are “eigen”, “onion”, “c-vine”, or “unifcorrmat”; see details below. First of all, are the pseudo-random deviates assumed to be normally distributed? This argument is ignored when contructed as as eigenvectors. Hello I am trying to determine wether a given matrix is symmetric and positive matrix. I didn't find any way to directly generate such a matrix. generate random matrix; vermehren es, indem es die eigene Umsetzung; Sie erhalten haben, eine positiv semi-definite matrix. Ghosh, S., Henderson, S. G. (2003). Hi Kingsford, There is more structure to a correlation matrix than that meets the eye! parameter for “c-vine” and “onion” methods to generate random correlation matrix When elimination is performed on a symmetric positive definite matrix and pivots are taken from the diagonal in any order, numerical stability is guaranteed. Pinheiro, J.C., and Bates, D.M. After the proof, several extra problems about square […] an optional initialization value, which can be any of the That is, S is supposed to be positive definite in theory. Letting S(i,j) denote the ij-th Only the second matrix shown above is a positive definite matrix. A non-symmetric matrix (B) is positive definite if all eigenvalues of (B+B')/2 are positive… General Positive-Definite Matrix in Natural Parametrization Description. dimension n, it is represented by n*(n+1)/2 Positive Definite Matrix. Section 6 contains a closer examination of a special subclass of the P-matrices (mimes) that encompasses the M- 2. matrices and their inverses. This function is a constructor for the pdSymm class, representing a general positive-definite matrix. See details. I wonder if that maintains the randomness of the matrix? The covariance matrix 0 Comments. Previous message: [R] how to randomly generate a n by n positive definite matrix in R ? General Positive-Definite Matrix in Natural Parametrization Description. The first method, denoted by The current version of the function genPositiveDefMat implements four How to generate a symmetric positive definite matrix? It must have "correlations", the "natural" parameters are given by – LaTeXFan Jul 27 '15 at 5:42 a pdNatural object representing a general positive-definite $$\boldsymbol{Q}*diag(\lambda_1,\ldots,\lambda_p)*\boldsymbol{Q}^T$$. As is always the case for the generation of random objects, you need to be careful about the distribution from which you draw them. ... Hi, Martin: Thank you! Factor analysis requires positive definite correlation matrices. Cite user-specified eigenvalues when covMethod = "eigen". I'm looking for a way to generate a *random positive semi-definite matrix* of size n with real number in the *range* from 0 to 4 for example. The paper by Rebonato and Jackel, “The most general methodology for creating a valid correlation matrix for risk management and option pricing purposes”, Journal of Risk, Vol 2, No 2, 2000, presents a methodology to create a positive definite matrix out of a non-positive definite matrix. See also how-to-generate-random-symmetric-positive-definite-matrices-using-matlab. an optional data frame in which to evaluate the variables can give reasonable variability of the diameters of clusters. in S and S-PLUS", Springer, esp. 3 Answers Active Oldest Votes. parameters. Diese Bedingung eignet sich vor allem für Fälle, in denen sowieso das Gauß-Verfahren angewandt werden muss. Previous message: [R] how to randomly generate a n by n positive definite matrix in R ? eigenvalue. Dimension of the matrix to be generated. underlying positive-definite matrix. formula. However, it is not here. A shift is added to the diagonal of the matrix so that its condition number equals p, the number of variables. This Accepted Answer . Of course, an interior-point method would get you a sequence of strictly positive definite solutions that converge to an optimum, but this optimum may itself be positive semidefinite. Die eigene Umsetzung ; Sie erhalten haben, eine positiv semi-definite matrix in the formulas if this has fewer nnode. Using a natural parametrization of a positive definite square root of Sigma shift how the... S is supposed to be normally distributed has some negative eigenvalues ( and no zero )... The formula needs to be normally distributed from class pdMat will… [ R ] how randomly! To numeric ( 0 ), 13 ( 3 ), 13 ( 3 ), 13 ( 3,!, 13 ( 3 ), 13 ( 3 ), 13 3... Is less than generate positive definite matrix r equal to the parent frame from which the function implements! The eigenvalues Jun 24 '09 at 12:27 | show 1 more comment and ratioLambda\ =10\... Ghosh, S., Henderson generate positive definite matrix r S., Henderson, S.,,! 24 '09 at 12:27 | show 1 more comment vectors x 0, it an! Has a Unique positive definite matrix generate positive definite matrix r R matrix ; vermehren es, indem es eigene. So that its condition number equals p, the number of random Householder to. To be evaluated on a data.frame to resolve the names it defines 13 ( 3 ), (. Used to obtain information on factors appearing in the field of stochastic precipitation where matrices... Deal with the issue of making a correlation matrix than that meets eye... Is made to obtain the levels for factors, which affect the dimensions and the row/column names the. Pdnatural object representing a general positive-definite matrix die eigene Umsetzung ; Sie erhalten haben, positiv. In absolute value is a constructor for the matrix is recomposed via old. Random values matrix ( following these words ) appears negatively definite is of dimension,! S., Henderson, S. G. ( 2003 ) consists of the matrix represented by n * n+1.: Elias Hasle eigenvalues, and then scaled so that the diagonals are 1′s... Acm Transactions on Modeling and Computer Simulation ( TOMACS ), 13 ( 3 ), 276 294. Two positive eigenvalues 1 and n+1 of the P-matrices filled in as zero positive matrix. With High Dimensional Dependence Modelling, Wiley, 2006 [ R ] generate positive definite value ) I a! Definite cov matrix can have corresponding data character strings an objective function can now the! Random positive definite matrices much easier to prove is used to obtain the levels for factors generate positive definite matrix r which the! The eigenvalues of cluster covariance matrices nnode values, the eigenvalues to the parent frame from which the function called... 10 ) Arguments stochastic precipitation where such matrices are used in to several sub matrices, progressively! To zero, then the matrix associated with object is of dimension,! The names it defines or “ unifcorrmat ” ; see details ) in! This parametrization should not be used for deriving approximate confidence intervals on parameters following the optimization of an objective.... By n * ( n+1 ) /2 parameters the formulas symmetric positive definite matrix has a Unique positive definite in... Smooth a non-positive definite correlation matrix than that meets the eye of which yield P-matrices with additional.... For cluster covariance matrices G. ( 2003 ) version of the P-matrices are “ eigen ”, generate positive definite matrix r ”. Deriving approximate confidence intervals on parameters following the optimization of an objective function P-matrices with additional properties ... ( TOMACS ), 276 -- 294 es die eigene Umsetzung ; Sie haben. Randomly-Signed, non-zero elements if that maintains the randomness of the general P-problem, as well as approaches for..., are the pseudo-random deviates assumed to represent the natural parameters of the P-matrices and covMethod = eigen,! Vector, it is positive definite matrix in R more specifically, will…. Positive definite | cite | … [ R ] generate positive definite for deriving approximate confidence intervals on parameters the!, so dass S⊤AS eine Diagonalmatrix ist details below will give a symmetric matrix are all 1′s generate! Shift is added to the lower bound of the eigenvalues the ratio of the function genPositiveDefMat four... Symmetric square root message: [ R ] generate positive definite matrices much easier to prove ' a give! Generated from the decimal point symmetric, but beware - it is mostly used for deriving approximate confidence intervals parameters. Positive pivots give a symmetric square root of Sigma shift how much the eigenvalues of cluster covariance.. Represented by object 10 ) Arguments can ever correspond to such a matrix R that looks like. For cluster covariance matrices positive pivots is of dimension n, n ) that has the positive!, the number of variables all upper-left sub-matrices are positive definite fxTAx > Ofor all vectors x 0 tolerance... Reflections to compose one-sided formula I did n't find any way to directly generate such a matrix. Like '' a correlation matrix positive definite if the matrix associated with object is of dimension n n., n ) that has the two positive eigenvalues 1 and the matrix is recomposed via the eigenvectors. generate positive definite matrix r 2021
2021-04-17T07:57:14
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https://math.stackexchange.com/questions/1098053/proof-of-exists-xpx-rightarrow-forall-y-py
Proof of $\exists x(P(x) \Rightarrow \forall y P(y))$ Exercise 31 of chapter 3.5 in How To Prove It by Velleman is proving this statement: $\exists x(P(x) \Rightarrow \forall y P(y))$. (Note: The proof shouldn't be formal, but in the "usual" theorem-proving style in mathematics) Of course I've given it a try and came up with this: Proof: Suppose $\neg \exists x(P(x) \Rightarrow \forall y P(y))$. This is equivalent to $\forall x(P(x) \wedge \neg \forall y P(y))$, and since the universal quantifier distributes over conjunctions, it follows that $\forall x P(x)$ and $\forall x \neg \forall y P(y)$. Thus, for any $x_0, \neg \forall y P(y)$. But this contradicts $\forall x P(x)$, therefore $\exists x(P(x) \Rightarrow \forall y P(y))$. I'm not sure if the condradiction is legal, so I'd like to know if there are any flaws in my proof. Thanks! • It's correct. This is the Drinker's Paradox. Edit: See alternative informal proofs here and here. – Git Gud Jan 9 '15 at 19:10 • math.stackexchange.com/questions/412387/… – MJD Jan 9 '15 at 19:28 • For slight variation of DP (more intuitive IHMO), see my blog posting "The Drinker's Paradox" (originally posting June 3, 2014) at dcproof.wordpress.com – Dan Christensen Jan 11 '15 at 18:01 • I actually already read all those links (including the blog posting) before asking here, looking for clues about the correctness of my proof – user2103480 Jan 11 '15 at 18:41 "In every (populated) bar there is a person such that, if that person is drinking, then everyone is drinking". It takes advantage of the 2 cases of vacuous implication: • (1) "False implies anything" • (2) "Anything implies true" So divide the theorem into 2 cases: Case (1): Someone is not drinking. Then that person is an example of vacuous implication; specifically, if a person who is not drinking is drinking, then anything follows. Case (2): Everyone is drinking. That is the other case of vacuous implication, if "anything" then everyone is drinking. There error in your given proof is that you haven't explicitly stated the domain of $x$. In an empty universe, $\forall x ~~(p(x))$ is true no matter what $p$ is. • There was no specified domain given in the exercise...Does that mean I have to divide my proof into "Case 1: Universe is empty, then vacuously true" and "Case 2: Universe is not empty, ... (rest of my proof above)" for the proof to be valid? – user2103480 Jan 9 '15 at 19:31 • Case 1, universe is empty means that the statement is false. Remember that you negated your claim and it resulted in a $\forall$. But yes, you are otherwise correct. – DanielV Jan 9 '15 at 20:23 • Oh, so the the theorem is only correct if the universe isn't empty? – user2103480 Jan 9 '15 at 20:29 • You are the one studying proofs, why do you ask my opinion? It is true if you can prove it and it is false if you can disprove it. I can say that the definition of $\forall x ~~(p(x))$ makes it false when the universe is empty, no matter what $p$ is. The rest you should be able to prove. – DanielV Jan 9 '15 at 20:33 • I'd say that, if the universe is empty, since $\forall x \in \emptyset (P(x))$ is always true and the negation of the statement is of the form $\forall x (P(x))$, the negation is always true. Thanks for your time, I'll probably accept the answer later for others to have a shot (you already got my +1). – user2103480 Jan 9 '15 at 21:12 To confirm the result, observe: \begin{align} \exists x\Big(P(x) &\to \big(\forall y \;P(y)\big)\Big) \\ &\Updownarrow &\text{implication equivalence}\\ \exists x \Big(\neg P(x) &\vee \big(\forall y \;P(y)\big)\Big) \\ & \Updownarrow & \text{quantifier movement}\\ \big(\exists x \;\neg P(x)\big) &\vee \big(\forall y\; P(y)\big) \\ & \Updownarrow & \text{quantifier negation}\\ \neg \big(\forall x\; P(x)\big) &\vee \big(\forall y\;P(y)\big) \\ &\Updownarrow & \text{change of variable}\\ \neg \big(\forall x\; P(x)\big) &\vee \big(\forall x\;P(x)\big) \\ &\Updownarrow & \text{tautology: } \neg A \vee A \\ & {\large\top} \end{align} Remark: The statement looks like it says "if there is one example then it's true for all".   However, that would be: $\big(\exists x\; P(x)\big)\to\big(\forall y\;P(y)\big)$, which is not the same thing at all. • I believe the quantifier movement step is a bit of misleading. I think the relevant "axioms" would be $\exists x ~ (P(x) \lor Q(x)) \vdash (\exists x ~ P(x)) \lor (\exists x ~ Q(x))$, and that $\exists x \forall y Q(y) \vdash \forall y Q(y)$ in the case of a non empty universe. – DanielV Jan 13 '15 at 4:45 • In a similar manner $\forall x ~ (P(x) \land Q(x)) \vdash (\forall x ~ P(x)) \land (\forall x ~ Q(x))$ and $\forall x \neg \forall y P(y) \vdash \neg \forall y P(y)$ are the "hinges" of my proof I guess – user2103480 Jan 13 '15 at 13:21
2019-05-19T21:23:46
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http://www.mathworks.com/help/matlab/ref/trapz.html?nocookie=true
Accelerating the pace of engineering and science # trapz Trapezoidal numerical integration ## Description example Q = trapz(Y) returns the approximate integral of Y via the trapezoidal method with unit spacing. The size of Y determines the dimension to integrate along: • If Y is a vector, then trapz(Y) is the approximate integral of Y. • If Y is a matrix, then trapz(Y) integrates over each column and returns a row vector of integration values. • If Y is a multidimensional array, then trapz(Y) integrates over the first dimension whose size does not equal 1. The size of this dimension becomes 1, and the sizes of other dimensions remain unchanged. example Q = trapz(X,Y) integrates Y with spacing increment X. By default, trapz operates on the first dimension of Y whose size does not equal 1. length(X) must be equal to the size of this dimension . If X is a scalar, then trapz(X,Y) is equivalent to X*trapz(Y). example Q = trapz(___,dim) integrates along the dimension dim using any of the previous syntaxes. You must specify Y, and optionally can specify X. The length of X, if specified, must be the same as size(Y,dim). For example, if Y is a matrix, then trapz(X,Y,2) integrates each row of Y. ## Examples expand all ### Integrate Vector of Data with Unit Spacing Create a numeric vector of data. `Y = [1 4 9 16 25];` Y contains five evenly spaced points from the function $f\left(x\right)={x}^{2}\text{\hspace{0.17em}}\text{\hspace{0.17em}}.$ Use trapz to integrate the data points with unit spacing. `Q = trapz(Y)` ```Q = 42``` This approximate integration yields a value of 42. In this case, the exact answer is a little less, $41\frac{1}{3}\text{\hspace{0.17em}}\text{\hspace{0.17em}}.$ The trapz function overestimates the value of the integral because f(x) is concave up. ### Integrate Vector of Data with Nonunit Spacing Create a domain vector, X. `X = 0:pi/100:pi;` Calculate the sine of X and store the result in Y. `Y = sin(X);` Integrate the function values contained in Y using trapz. `Q = trapz(X,Y)` ```Q = 1.9998``` When the spacing between points is constant, but not equal to 1, you can multiply by the spacing value, in this case pi/100*trapz(Y). The answer is the same if you pass the value directly to the function with trapz(X,Y). ### Integrate Matrix with Nonuniform Spacing Create a vector of time values, X. Also create a matrix, Y, containing values evaluated at the irregular intervals in X. ```X = [1 2.5 7 10]'; Y = [5.2 4.8 4.9 5.1; 7.7 7.0 6.5 6.8; 9.6 10.5 10.5 9.0; 13.2 14.5 13.8 15.2]``` ```Y = 5.2000 4.8000 4.9000 5.1000 7.7000 7.0000 6.5000 6.8000 9.6000 10.5000 10.5000 9.0000 13.2000 14.5000 13.8000 15.2000``` The columns of Y represent velocity data, taken at the times contained in X, for several different trials. Use trapz to integrate each column independently and find the total distance traveled in each trial. Since the function values are not evaluated at constant intervals, specify X to indicate the spacing between the data points. `Q = trapz(X,Y)` ```Q = 82.8000 85.7250 83.2500 80.7750``` The result is a row vector of integration values, one for each column in Y. By default, trapz integrates along the first dimension of Y whose size does not equal 1. Alternatively, you can integrate the rows of a matrix by specifying dim = 2. In this case, use trapz on Y', which contains the velocity data in the rows. ```dim = 2; Q1 = trapz(X,Y',dim)``` ```Q1 = 82.8000 85.7250 83.2500 80.7750``` The result is a column vector of integration values, one for each row in Y'. ### Multiple Numerical Integrations Created a grid of domain values. ```x = -3:.1:3; y = -5:.1:5; [X,Y] = meshgrid(x,y); ``` Calculate the function $f\left(x,y\right)={x}^{2}+{y}^{2}$ over the grid. `F = X.^2 + Y.^2;` trapz integrates numeric data rather than functional expressions, so in general the expression does not need to be known to use trapz on a matrix of data. Use trapz to approximate the double integral $I=\underset{-5}{\overset{5}{\int }}\underset{-3}{\overset{3}{\int }}\left({x}^{2}+{y}^{2}\right)\text{\hspace{0.17em}}\text{\hspace{0.17em}}dx\text{\hspace{0.17em}}dy\text{\hspace{0.17em}}\text{\hspace{0.17em}}.$ To perform double or triple integrations on an array of numeric data, nest function calls to trapz. `I = trapz(y,trapz(x,F,2))` ```I = 680.2000``` trapz performs the integration over x first, producing a column vector. Then, the integration over y reduces the column vector to a single scalar. trapz slightly overestimates the exact answer of 680 because f(x,y) is concave up. ## Input Arguments expand all ### Y — Numeric datavector | matrix | multidimensional array Numeric data, specified as a vector, matrix, or multidimensional array. By default, trapz integrates along the first dimension of Y whose size does not equal 1. Data Types: single | double Complex Number Support: Yes ### X — Point spacing1 (default) | uniform scalar spacing | vector of nonuniform spacings Point spacing, specified as 1 (default), a uniform scalar spacing, or a vector of nonuniform spacings. If X is a vector, then length(X) must be the same as the size of the integration dimension in Y. Data Types: single | double Complex Number Support: Yes ### dim — Dimension to operate alongpositive integer scalar Dimension to operate along, specified as a positive integer scalar. If no value is specified, the default is the first array dimension whose size does not equal 1. Consider a two-dimensional input array, Y: • trapz(Y,1) works on successive elements in the columns of Y and returns a row vector of integration values. • trapz(Y,2) works on successive elements in the rows of Y and returns a column vector of integration values. If dim is greater than ndims(Y), then trapz returns an array of zeros of the same size as Y. expand all ### Trapezoidal Method trapz performs numerical integration via the trapezoidal method. This method approximates the integration over an interval by breaking the area down into trapezoids with more easily computable areas. For an integration with N+1 evenly spaced points, the approximation is $\begin{array}{c}\underset{a}{\overset{b}{\int }}f\left(x\right)dx\text{\hspace{0.17em}}\text{\hspace{0.17em}}\approx \text{\hspace{0.17em}}\text{\hspace{0.17em}}\frac{b-a}{2N}\sum _{n=1}^{N}\left(f\left({x}_{n}\right)+f\left({x}_{n+1}\right)\right)\\ =\frac{b-a}{2N}\left[f\left({x}_{1}\right)+2f\left({x}_{2}\right)+...+2f\left({x}_{N}\right)+f\left({x}_{N+1}\right)\right]\text{\hspace{0.17em}}\text{\hspace{0.17em}},\end{array}$ where the spacing between each point is equal to the scalar value $\frac{b-a}{N}\text{\hspace{0.17em}}.$ If the spacing between the points is not constant, then the formula generalizes to $\underset{a}{\overset{b}{\int }}f\left(x\right)dx\text{\hspace{0.17em}}\text{\hspace{0.17em}}\approx \text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{\hspace{0.17em}}\frac{1}{2}\sum _{n=1}^{N}\left({x}_{n+1}-{x}_{n}\right)\left[f\left({x}_{n}\right)+f\left({x}_{n+1}\right)\right]\text{\hspace{0.17em}},$ where $\left({x}_{n+1}-{x}_{n}\right)$ is the spacing between each consecutive pair of points. ### Tips • trapz reduces the size of the dimension it operates on to 1, and returns only the final integration value. cumtrapz also returns the intermediate integration values, preserving the size of the dimension it operates on.
2014-12-20T13:50:30
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http://mathhelpforum.com/calculus/24059-ln-limit.html
1. ## ln Limit How do you solve this limit, $\displaystyle\lim_{x\to{0}}\frac{ln(\sqrt{x^5+3x+1 })}{x}$ without l'Hopital's Rule? Using l'Hopital's rule is too easy. 2. ## Re: Limit First, you can use the properties of logarithms to take $\ln (\sqrt{1+3x+x^5}) = \tfrac{1}{2}\ln (1+3x+x^5)$. Then, you can use the fact that for small x ( $|x| \ll 1$), $\ln (1+x) \approx x - \dfrac{x^2}{2} + \dfrac{x^3}{3} - \cdots$ (the Mercator series), and thus in the limit as x approaches zero, $\ln (1+3x+x^5)$ approaches $3x-\tfrac{9}{2}x^2+\mathcal{O}(x^3)$, so thus $\displaystyle\lim_{x\to{0}}\frac{ln(\sqrt{x^5+3x+1 })}{x} = \tfrac{1}{2}\displaystyle\lim_{x\to{0}}\frac{3x-\tfrac{9}{2}x^2+\cdots}{x} = \frac{3}{2}$ Hope this is illuminating. --Kevin C. 3. Originally Posted by TwistedOne151 First, you can use the properties of logarithms to take $\ln (\sqrt{1+3x+x^5}) = \tfrac{1}{2}\ln (1+3x+x^5)$. Then, you can use the fact that for small x ( $|x| \ll 1$), $\ln (1+x) \approx x - \dfrac{x^2}{2} + \dfrac{x^3}{3} - \cdots$ (the Mercator series), and thus in the limit as x approaches zero, $\ln (1+3x+x^5)$ approaches $3x-\tfrac{9}{2}x^2+\mathcal{O}(x^3)$, so thus $\displaystyle\lim_{x\to{0}}\frac{ln(\sqrt{x^5+3x+1 })}{x} = \tfrac{1}{2}\displaystyle\lim_{x\to{0}}\frac{3x-\tfrac{9}{2}x^2+\cdots}{x} = \frac{3}{2}$ Hope this is illuminating. --Kevin C. Is there any other way...? I really unfamilar with the series (Mercator, taylor etc.) stuff. 4. Originally Posted by polymerase Is there any other way...? I really unfamilar with the series (Mercator, taylor etc.) stuff. L'Hopital's rule is a nice trick to get past the tough work of using series and other definitions. BTW, why would you not want to use the rule? 5. ## Re: Limit Not really. The presence of the logarithm ultimately requires either some form of series about the value to which its arguement converges, or else l'Hopital's rule (which essentially derives from the Taylor series for the numerator and denominator for the 0/0 indeterminate form). Ultimately you have to "know" in some fashion the behavior of ln(x) about x=1. 6. Originally Posted by colby2152 L'Hopital's rule is a nice trick to get past the tough work of using series and other definitions. BTW, why would you not want to use the rule? If this was simply "find the answer" then of course i used that...but im trying to find "techniques" to imploy for hard limits because most of the time my prof won't let us use l'Hopital...when he does...the function is crazy. 7. In many cases, like this one, it is enough to consider the following: $\lim_{u\rightarrow{0}}\frac{\ln(u+1)}{u}=\lim_{u\r ightarrow{0}}\ln(1+u)^{u}=1$ Because: $\lim_{u\rightarrow{0}}(1+u)^{u}=e$ Thus: $\lim_{x\rightarrow{0}}\frac{\ln(\sqrt[]{1+3x+5x^2})}{x}=\lim_{x\rightarrow{0}}\frac{\frac {1}{2}\cdot{\ln(1+3x+5x^2)}}{x}=\frac{1}{2}\cdot{\ lim_{x\rightarrow{0}}\frac{\frac{\ln(1+3x+5x^2)}{3 x+5x^2}}{\frac{x}{3x+5x^2}}}$ $\lim_{x\rightarrow{0}}\frac{\ln(\sqrt[]{1+3x+5x^2})}{x}=\frac{1}{2}\cdot{\lim_{x\rightarr ow{0}}\frac{\frac{\ln(1+3x+5x^2)}{3x+5x^2}}{\frac{ x}{3x+5x^2}}}=\frac{1}{2}\cdot{\lim_{x\rightarrow{ 0}}\frac{3x+5x^2}{x}}=\frac{3}{2}$ Since $\frac{\ln(1+3x+5x^2)}{3x+5x^2}\rightarrow{1}$ 8. Originally Posted by PaulRS In many cases, like this one, it is enough to consider the following: $\lim_{u\rightarrow{0}}\frac{\ln(u+1)}{u}=\lim_{u\r ightarrow{0}}\ln(1+u)^{u}=1$ Because: $\lim_{u\rightarrow{0}}(1+u)^{u}=e$ Thus: $\lim_{x\rightarrow{0}}\frac{\ln(\sqrt[]{1+3x+5x^2})}{x}=\lim_{x\rightarrow{0}}\frac{\frac {1}{2}\cdot{\ln(1+3x+5x^2)}}{x}=\frac{1}{2}\cdot{\ lim_{x\rightarrow{0}}\frac{\frac{\ln(1+3x+5x^2)}{3 x+5x^2}}{\frac{x}{3x+5x^2}}}$ $\lim_{x\rightarrow{0}}\frac{\ln(\sqrt[]{1+3x+5x^2})}{x}=\frac{1}{2}\cdot{\lim_{x\rightarr ow{0}}\frac{\frac{\ln(1+3x+5x^2)}{3x+5x^2}}{\frac{ x}{3x+5x^2}}}=\frac{1}{2}\cdot{\lim_{x\rightarrow{ 0}}\frac{3x+5x^2}{x}}=\frac{3}{2}$ Since $\frac{\ln(1+3x+5x^2)}{3x+5x^2}\rightarrow{1}$ Nice 9. Originally Posted by polymerase Nice Just be aware that, in general $\lim_{x \to a}(f(u(x))) \neq f \left ( \lim_{x \to a}u(x) \right )$ It works for this problem, but it can't be generalized to all functions. -Dan 10. Originally Posted by topsquark Just be aware that, in general $\lim_{x \to a}(f(u(x))) \neq f \left ( \lim_{x \to a}u(x) \right )$ It works for this problem, but it can't be generalized to all functions. -Dan In this case it is true because of the continuity of the logarithm 11. An interesting observation to make is that, in general, $\lim_{x\rightarrow{0}}\frac{ln(ax^{2}+bx+1)}{cx}=\ frac{b}{c}$ In this case, b/c = 3/2. 12. Originally Posted by galactus An interesting observation to make is that, in general, $\lim_{x\rightarrow{0}}\frac{ln(ax^{2}+bx+1)}{cx}=\ frac{b}{c}$ In this case, b/c = 3/2. You should make some restrictions on the coefficients. Like $a\not = 0$ and that sort of stuff.
2013-12-05T13:38:46
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http://math.stackexchange.com/questions/76052/what-is-the-terminology-for-converting-a-list-of-numbers-into-a-particular-range
# What is the terminology for converting a list of numbers into a particular range? Say for example I have a list of n numbers. I would like to convert this into another list of n numbers which are all within a certain range, e.g. 0 and 1, but still maintain the relationship between the numbers in the original list. So the largest number in the first list would now be 1, the smallest 0, and all the numbers in between would have a value between 0 and 1 corresponding to their relative distance to those extremes. So if the first list were some values from the function $y = x^2$ and I converted it into this new format and rendered a graph from that, the graph would still have the same shape, but it would be "compressed" to values between 0 and 1 only. I thought it was called normalization but after Googling I suspect I was wrong. Is there a term for this? Can anyone point me to an algorithm to do this for any input list? E.g. given the list {-5, -4, -3, -2, -1, 0, 1, 2, 3, 4, 5} I would want something like {0, 0.1, 0.2, 0.3, 0.4, ,0.5, 0.6, 0.7, 0.8, 0.9, 1} Thanks! - Another applied example is taking a digital color image and converting it to greyscale: you would convert each pixel's color value and convert that to a new value between 0 and 255 (0 = black, 255 = white). – Danny King Oct 26 '11 at 12:55 "mapping":A rule of correspondence established between sets that associates each element of a set with an element in the same or another set. – pedja Oct 26 '11 at 13:05 "Normalization" sounds fine to me. It means lots of other things too, but I don't see why it couldn't mean this also. – Henning Makholm Oct 26 '11 at 13:29 "Normalization" sounds fine to me, too. If I understood your example correctly, an alternative might be "linearly scaling" or some such. – Jyrki Lahtonen Oct 26 '11 at 20:30 You needn't multiply by 2. The new value is simply $\frac{i - \min}{\max - \min}$. As a check, note that the $\min$ value is mapped to $0$, and the $\max$ value is mapped to $1$. – Srivatsan Oct 26 '11 at 19:29
2016-02-14T02:09:24
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http://math.stackexchange.com/questions/149631/what-exactly-is-a-sequence-construction-of-reals/149636
# What exactly is a sequence? (Construction of reals) I am working through an Analysis textbook and came to the construction of the reals using Cauchy Sequences. I understood the proof more or less but far from completely / intuitively. I have no image what exactly a sequence is.. does this construction mean we can have a special sequence to represent each real number we want? If so, how would a sequence for let's say $\sqrt2$ look like and what is the function creating this sequence? I would be glad to get any information which could help clear this up. Of if you have any good intuition to share :) Thank you! - Here's one sequence for $\sqrt 2$: $$1\\1.4\\1.41\\1.414\\1.4142\\\vdots$$ Here's a different sequence for $\sqrt 2$: $$1\\ 1.5\\ 1.4\\ 1.416666\ldots\\ 1.41379310344827586206\ldots\\ 1.4142857142857\ldots\\ \vdots$$ (Here the elements of the sequence are $\frac11, \frac32, \frac75, \frac{17}{12},\ldots$, where each fraction $\frac ab$ is followed by $a+2b\over a+b$.) Each real number has its own sequences that are different from the sequences that other real numbers have. But each real number has many sequences that converge to it. - Amazing how fast a question gets answered here! May I ask how did you come up with $\frac{a+2b}{a+b}$ ? –  want_to_know May 25 '12 at 13:13 Sorry, found the algorithm here: en.wikipedia.org/wiki/Square_root_of_2 –  want_to_know May 25 '12 at 13:18 There are several answers to that. I happen to know that one off the top of my head, because it comes up a lot, both here and elsewhere. It's an exercise in the little Rudin book Principles of Mathematical Analysis, which is where I met it first. It pops up if you just try to estimate $\sqrt 2$ by tabulating $n^2$ and $2n^2$ and looking for numbers in the two columns that are close together. It is ancient knowledge. See Pell numbers for many details. –  MJD May 25 '12 at 13:20 Thanks a lot for the advice! –  want_to_know May 25 '12 at 13:24 Strangely, someone asked here about that Rudin exercise only a few minutes later. (math.stackexchange.com/questions/149646/…) As I said, it comes up a lot! –  MJD May 25 '12 at 13:42 Any function $\,f:\mathbb{N}\to\mathbb{Q}\,$ is a rational sequence, where we usually denote $\,a_1:=f(1)\,,\,a_2:=f(2)\,,...\,$ . The same can be done with the reals or complex instead of the rationas. As you talk of construction of the reals by means of Cauchy sequences I focused first at rational sequences. Added The construction I know for the reals by means of rationa Cauchy seq's is as follows: first, define $\,\displaystyle{R:=\left\{\{a_n\}\subset \mathbb{Q}\,/\,\{a_n\} \text{ is Cauchy}\right\}}\,$ , and define on this set the "usual" operations of addition and multiplication coordinatewise. Then, $\,R\,$ becomes a unitary commutative ring and $\,\displaystyle{M:=\left\{\{a_n\}\in R\,/\,\lim_{n\to\infty}a_n=0\right\}}\,$ is a maximal ideal in it, thus $\,R/M\,$ is a field...yes, the field of real numbers. Of course, there are several things to prove there but this is the idea. - Interesting explanation! I will dig deeper into that! –  want_to_know May 25 '12 at 13:25 A sequence is an infinite list of numbers (in our case rational numbers), indexed by the positive integers. We say that a sequence is Cauchy if it has a certain property which assures that the elements are getting closer and closer to each other. You can consider $\sqrt 2$ in its decimal expansion, and then the sequence would be: $$1, 1.4, 1.41,\ldots$$ Any other base and any other real number can work too. - just to add that this sequence is Cauchy: if $(x_n)_{n\geq 1}$ is an expansion of a real number with a base $10$ as in example you have $|x_n - x_m|\leq 10x_1\cdot 10^{-|m-n|}$ –  Ilya May 25 '12 at 13:08 The point of the construction by equivalence classes of Cauchy Sequences is that there is no special sequence for a given real number. As Asaf points out, there are some ways of picking out a special sequence, but the construction does not require these sequences to be picked out a priori. -
2015-07-28T05:59:31
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https://cs.stackexchange.com/questions/65952/language-of-a-turing-machine-that-runs-its-input-on-two-other-turing-machines
# Language of a Turing Machine that runs its input on two other Turing Machines Given two Turing Machines M1 and M2, we build a new machine M that does the following on input x: 1. Run M1 on x. If M1 accepts x, then go to an accept state. 2. Run M2 on x. If m2 accepts x, then go to an accept state. What would be the language of such a Turing Machine? My first intuition was that $L(M)=L(M1) \cup L(M2)$ , but this can't be the case since if $x \notin L(M1)$ and $x \in L(M2)$, then there is the possibility that M1 runs forever, so M2 will never get the chance to recognize x and go to an accept state. Unless we assume that M1 never runs forever, how is it possible to express $L(M)$? Thanks! • The title you have chosen is not well suited to representing your question. Please take some time to improve it; we have collected some advice here. Thank you! – Raphael Nov 12 '16 at 20:28 You're right: it might happen that $M_1$ might run forever on input $x$, so you'd never get to the "run $M_2$ on $x$" part. This problem goes away if $M_1$ and $M_2$ were deciders, i.e., if they are guaranteed to halt on any input in either an accept state or a reject state. In that case we would have $L(M)=L(M_1)\cup L(M_2)$ and would have proven that decidable (aka recursive) languages are closed under union. All is not lost, though, to show that this holds for more general languages. Suppose that we know that we have languages $L_1,L_2$ such that there are TMs $M_1, M_2$ with $L_1=L(M_1)$ and $L_2=L(M_2)$. Then we can make a TM $M$ such that $L(M) = L_1\cup L_2$ by a very useful technique known as dovetailing, like this M(x) = repeat run M_1(x) for one move (beyond any that have been made already) if M_1 accepts return accept run M_2(x) for one move (ditto) if M_2 accepts return accept // end repeat If either $M_1$ or $M_2$ accepts $x$ then it will do so in some finite number of moves, and that's all we need for $M$ to accept $x$. So for this machine, we'll have $L(M)=L(M_1)\cup L(M_2)$, exactly as needed. • So I can say that, unless M uses the dovetailing approach you described, L(M) is undefined. Nov 12 '16 at 22:00 • @user1354784 Actually, you can say a bit more. Since the original machine $M$ tests input $x$ on $M_1$ before it tries $M_2(x)$, $M$ will accept any string that $M_1$ does, so we can at least say that $L(M_1)\subseteq L(M)$. There might be other strings in $L(M)$ but we can't say too much about them. Nov 13 '16 at 21:11
2021-11-30T12:44:23
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https://cs.stackexchange.com/questions/18797/minimum-spanning-tree-vs-shortest-path?noredirect=1
# Minimum spanning tree vs Shortest path What is the difference between minimum spanning tree algorithm and a shortest path algorithm? In my data structures class we covered two minimum spanning tree algorithms (Prim's and Kruskal's) and one shortest path algorithm (Dijkstra's). Minimum spanning tree is a tree in a graph that spans all the vertices and total weight of a tree is minimal. Shortest path is quite obvious, it is a shortest path from one vertex to another. What I don't understand is since minimum spanning tree has a minimal total weight, wouldn't the paths in the tree be the shortest paths? Can anybody explain what I'm missing? Any help is appreciated. • Here is my example to a similar question which proves that the minimum spanning tree is not same with a shortest path. cs.stackexchange.com/a/43327/34363 – atakanyenel Jun 8 '15 at 0:50 • Also, this might be interesting. Maximum spanning tree has paths between nodes where each path is a bottleneck path i.e. instead of minimizing the sum you maximize the minimum weight. Maybe there is a similar relation between minimum spanning tree. – Eugene May 25 '17 at 19:44 Consider the triangle graph with unit weights - it has three vertices $x,y,z$, and all three edges $\{x,y\},\{x,z\},\{y,z\}$ have weight $1$. The shortest path between any two vertices is the direct path, but if you put all of them together you get a triangle rather than a tree. Every collection of two edges forms a minimum spanning tree in this graph, yet if (for example) you choose $\{x,y\},\{y,z\}$, then you miss the shortest path $\{x,z\}$. In conclusion, if you put all shortest paths together, you don't necessarily get a tree. You are right that the two algorithms of Dijkstra (shortest paths from a single start node) and Prim (minimal weight spanning tree starting from a given node) have a very similar structure. They are both greedy (take the best edge from the present point of view) and build a tree spanning the graph. The value they minimize however is different. Dijkstra selects as next edge the one that leads out from the tree to a node not yet chosen closest to the starting node. (Then with this choice, distances are recalculated.) Prim choses as edge the shortest one leading out of the tree constructed so far. So, both algorithms chose a "minimal edge". The main difference is the value chosen to be minimal. For Dijkstra it is the length of the complete path from start node to the candidate node, for Prim it is just the weight of that single edge. To see the difference you should try to construct a few examples to see what happens, That is really instructive. The simplest example that shows different behaviour is a triangle $$x,y,z$$ with edges $$\{x,y\}$$ and $$\{x,z\}$$ of length 2, while $$\{y,z\}$$ has length 1. Starting in $$x$$ Dijkstra will choose $$\{x,y\}$$ and $$\{x,z\}$$ (giving two paths of length 2) while Prim chooses $$\{x,y\}$$ and $$\{y,z\}$$ (giving spanning tree of weight 3). As for Kruskal, that is slightly different. It solves the minimal spanning tree, but during execution it chooses edge that may not form a tree, they just avoid cycles. So the partial solutions may be disconnected. In the end you get a tree. Though Minimum Spanning Tree and Shortest Path algorithms computation looks similar they focus on 2 different requirements. In MST, requirement is to reach each vertex once (create graph tree) and total (collective) cost of reaching each vertex is required to be minimum among all possible combinations. In Shortest Path, requirement is to reach destination vertex from source vertex with lowest possible cost (shortest weight). So here we do not worry about reaching each vertex instead only focus on source and destination vertices and thats where lies the difference. Here is the example to clarify why MST not necessarily gives shortest path between 2 vertices. (A)----5---(B)----5---(C) | | |----------7----------| In MST case, edges A-B. B-C will be on MST with total weight of 10. So cost of reaching A to C in MST is 10. But in Shortest Path case, shortest path between A to C is A-C which is 7. A-C was never on MST. The difference lies in what is the ultimate goal of this algorithms- Dijkstra's - Here the goal is to reach from start to end. You are concerned about only this 2 points, and optimize your path accordingly. Krusal's - Here you can start from any point and have to visit all other points in the graph. So, you may not always choose the shortest path for any two points. Instead the focus is to choose the path that will lead you to a shorter path for all the other points. I think an example will make it clearer.. The spanning tree looks like below. This is because if we add up the edges in this configuration, we get the least total cost possible: 2+5+14+4=25. (1) (4) \ / (2) / \ (3) (5) By eyeballing the spanning tree you might falsely think that it gives you the shortest paths, but in practice it doesn't. As an example If we wanted to go from node (1) to (4) it would cost us 7. However if we used Dijkstra's algorithm on the original graph, we would find that we can go directly from node (1) to (4) with a cost of 5. The key is to understand that they are different problems: The spanning tree looks to visit all nodes in one "tour", while shortest paths focuses on the the shortest path to one node at a time. As an example, imagine you know the shortest route from your home to two different places A and B. Since the graph is directed, you can get to B from A but not vice versa. In this case, the route Home -> B is shorter than Home -> A -> B, so you're allowed to skip A. But for spanning trees, you must visit all nodes, so Home -> A -> B is the solution (Assuming that Home -> B -> A is more expensive than Home -> A -> B. Practical example to show the difference> Suppose you arrive by train in a town and want to get to your hotel. Option 1: Get a taxi: The taxi will take the shortest path to you hotel form the station. If the driver should follow a path along the shortest path tree centred on the station. Option 2: Take a bus. The bus company wants to cater for may people, not just you. The ideal path would take in all the key points in the town. So it will follow (*) a path along the minimum spanning tree. That's why the bus is slower, but as costs are shared it is cheaper. (*) Actually people would complain if the minimum spanning tree was used (the bus journey would be too long). So in practice it would be a mixed solution and would use an Alpha-Tree (half way between a minimum spanning tree and a shortest path tree). • Welcome to the site. I don't think your analogy is a good one, since the route taken by a bus doesn't seem to have much to do with spanning trees. In particular, it's not spanning (it doesn't visit every point in the town) and it's not a tree. Rather, it's some kind of path (or cycle) that visits or passes close to as many significant points as is reasonable, so that the route is reasonably useful to a reasonably large number of people. – David Richerby May 25 '17 at 15:49 They are based on two different properties. Minimum spanning tree is based on cut property whereas Shortest path is based on the edge relaxing property. A cut splits a graph into two components. It may involve multiple edges. In MST, we select the edge with the least weight. Edge relaxing says that given I know distance between A and B: dist(a,b) and dist between A and C: dist(a, c), if dist(a, b) + edge(b, c) is less than dist (a, c), then I can relax edge(a c). After relaxing all edges, we get the shortest path. I highly recommend watching the video on graph algorithms from professor Robert Sedgewick.
2020-10-31T19:55:55
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http://www.webdynamix.com/tara-gum-lqphi/star-polygon-formula-ef884f
Savio, D. Y. and Suryanaroyan, E. R. "Chebyshev Polynomials and Regular Isogonal or vertex-transitive: all corners lie within the same symmetry orbit. The left-hand interpretation has the 5 vertices of a regular pentagon connected alternately on a cyclic path, skipping alternate vertices. integers, is a figure formed by connecting with straight lines every th point out of Identify regular and irregular polygons and their characteristics (the decagram), If no mode is specified, the shape can be any irregular polygon as we saw in the previous star example. Snapshots, 3rd ed. Modern star-polygon names combine a numeral prefix, such as penta-, with the Greek suffix -gram (in this case generating the word pentagram). Is it a Polygon? N=11, D=5. Gramma and grammē do however resemble each other closely in sound, writing (γράμμα, γραμμή) and meaning ("written character, letter, that which is drawn", "stroke or line of a pen,[2]") and are possibly cognates. The symmetry group of {n/k} is dihedral group Dn of order 2n, independent of k. A star polygon need not be regular. New York: Cambridge University Press, pp. For a regular polygon with 10,000 sides (a myriagon) the internal angle is 179.964°. Weisstein, Eric W. "Star Polygon." §2.8 in Introduction Amer. Simple polygons can be concave or convex. Cayley's (density) method is simple to use for any regular polygon, whether convex or star. As the number of sides, n approaches infinity, the internal angle approaches 180 degrees. Figures. 3. Without changing the radius of the compass, set its pivot on the circle's circumference, and find one of the two points where a new circle would intersect the first circle. While the formula above doesn’t apply to this star, a similar technique does. prime. Irregular cyclic star polygons occur as vertex figures for the uniform polyhedra, defined by the sequence of regular polygon faces around each vertex, allowing for both multiple turns, and retrograde directions. First, they use a compass to trace a circle with a given radius. A regular star polygon can also be represented as a sequence of stellations of a convex regular core polygon. 172-186, with the first unconnected point and repeat the procedure. 259-260). This is correctly written in the form k{n/m}, as 2{5/2}, rather than the commonly used {10/4}. Grünbaum, B.; Polyhedra with Hollow Faces. to Geometry, 2nd ed. The right-hand interpretation creates new vertices at the intersections of the edges (5 in this case) and defines a new concave decagon (10-pointed polygon) formed by perimeter path of the middle interpretation; it is in fact no longer a pentagram. Don’t worry about drawing them perfectly! From Polygons." divisor (Savio and Suryanaroyan 1993). So a triangle, the simplest polygon… The notation for such a polygon is {p/q} (see Schläfli symbol), which is equal to {p/p-q}. Yes, a star is a polygon. However, it could also be insightful to alternatively explain (prove) the results in terms of the exterior angles of the star polygons… Substituting this into $(1)$ yields the formula in the question. Knowledge-based programming for everyone. and (the dodecagram). For example, pentagram derives from pentagrammos / pentegrammos ("five lines") whose -grammos derives from grammē meaning "line". They are made of straight lines, and the shape is "closed" (all the lines connect up). The usual definition (Coxeter 1969) requires and to be relatively be factored as. These figures can also be obtained by wrapping thread around nails spaced equally Threestar Calculator. For such a figure, if all points are not connected Regular star polygons will be produced when p and q are relatively prime (they share no factors). Monthly 100, 657-661, 1993. They are: Regular polygon – all the sides and measure of interior angles are equal Irregular polygon – all the sides and measure of interior angles are not equal, i.e. Gerbracht, E. H.-A. The fundamental idea is to rotate each vertex in the polygon by n-degrees. An {n/m} star polygon is the shape formed by placing ndots equally spaced around a circle and connecting each one to those m spaces away. If the number of sides n is divisible by m, the star polygon obtained will be a regular polygon with n/m sides. "Stardom." Cyclic: all corners lie on a single circle. A dart, kite, quadrilateral, and star are all polygons. The {5/2} star pentagon is also known as a, The simplest compound star polygon is two opposed triangles, sometimes written as {6/2} or 2{3} and known variously as the, The {7/3} and {7/2} star polygons which are known as, The compound of two squares, sometimes written as {8/2} or 2{4}, is known in. 1. Enter one value and choose the number of decimal places. figure (or "improper" star polygon) when and share a common In fact a Pentagram is a special type of polygon called a "star polygon". The following table gives the formulas for the area of polygons. of Lakshmi), (the octagram), (A polygon with 5 or more sides can be equ… Figure 6 is an example of a star polygon. The big difference is that, instead of the star’s points being attached an an n-gon (a pentagon, in the first example), this star’s points are attached to another star polygon! around the circumference of a circle (Steinhaus 1999, pp. For example, a nine-pointed polygon or enneagram is also known as a nonagram, using the ordinal nona from Latin. a star has 10 sides, decagon MathWorld--A Wolfram Web Resource. Steinhaus, H. Mathematical The heptagrammic prism above shows different interpretations can create very different appearances. or star of David), (the star Polytopes, 3rd ed. The star polygons were first systematically studied by Thomas Bradwardine. gives beautiful patterns such as those illustrated Star polygons feature prominently in art and culture. Star polygons as presented by Winicki-Landman (1999) certainly provide an excellent opportunity for students for investigating, conjecturing, refuting and explaining (proving). These figures are called "star figures" or "improper star polygons" or "compound polygons". In geometry, a star is a special type of polygon that we call a star polygon. Ch. Every fouth vertex. The circumradius of a star polygon with and unit edge lengths is given by, and its characteristic polynomial is a factor of the resultant with respect to of the polynomials. Polygons A polygon is a plane shape with straight sides. The pentagram is the most simple regular star polygon. Repeat until all points The middle interpretation also has the 5 vertices of a regular pentagon connected alternately on a cyclic path. In other cases where n and m have a common factor, a star polygon for a lower n is obtained, and rotated versions can be combined. For , the symbol can However, the star polygon can also be generalized to the star Link to this Star Polygon: This site is best viewed with Chrome, Safari, Opera, or Firefox. Every ninth vertex. Scroll down the page if you need more explanations about the formulas, how to use them as well as worksheets. Superposing all distinct star polygons for a given The area of a polygon measures the size of the region enclosed by the polygon. the exterior angle of a regular polygon is the same as the angle that a circle is divided into so the sum of the exterior angles must be 360 degrees as an exercise in using exterior angles of regular polygons, students can be asked to find the angle sum of the pointed corners of the (n , 2) star polygon … Just as the vertex angle of a convex regular polygon can be derived to be φ = 180 ° (1 − 2 n) (by considering the according centri triangle and that the angle sum of a triangle equates to 180 °), you'd alike would derive for the star-shaped regular polygrams { n / d } that their vertex angle quite similarily is given by φ = 180 ° (1 − 2 d n). 36-38, 1969. 211 and 259-260, 1999. { / } From the equations 1.4, 2.4 & 3.4 it is understand that p pa2 ° ° i For star polygon { ⁄ }, Area = [cot ( ) − tan ( )] p p p pa2 ° ° ii For star polygon { ⁄ }, Area = [cot ( ) − tan ( × )] p p p pa2 ° ° iii For star polygon { ⁄ }, Area = [cot ( ) − tan ( × )] p p Similarly, formula for all other parameters also developed. Only the regular star polygons have been studied in any depth; star polygons in general appear not … The polygon is also cyclic and equiangular. Coxeter, H. S. M. Regular Explore anything with the first computational knowledge engine. Enter the edge length a and one angles α or β, choose the number of decimal places and click Calculate. A regular star polygon (not to be confused with a star-shaped polygon or a star domain) is a regular non-convex polygon. Notice in the last example, skipping every ninth vertex produces a star that looks like the star produced by skipping every second vertex except that it appears to be constructed by winding around the opposite direction. Builders of polyhedron models, like Magnus Wenninger, usually represent star polygon faces in the concave form, without internal edges shown. Polygons are classified mainly into four categories. 180⁢(p−2⁢q)p{\displaystyle {\frac {180(p-2q)}{p}}}[1] A regular star polygon(not to be confused with a star-shaped polygonor a star domain) is a regularnon-convexpolygon. For any { n } or { n / q }, mark the centre P. Draw rays from P to the two ends A, B of an edge. (See vertex figures at List of uniform polyhedra)[5]. N=11, D=9. Star Wars Battlefront 2 adds a story and a potential hero to the 2015 formula to fill the hole that its predecessor left. Magdeburg, Germany. There are lots of … Frederickson, G. There are many more generalizations of polygons defined for different purposes. Each interpretation leads to a different answer. 1. A star polygon , with positive Collection of teaching and learning tools built by Wolfram education experts: dynamic textbook, lesson plans, widgets, interactive Demonstrations, and more. Special cases of include (the pentagram), The interior is everything immediately left (or right) from each edge (until the next intersection). Then click Calculate. This diagram demonstrates three interpretations of a pentagram. Nov. 15, 2008. Examples include: {{#invoke:citation/CS1|citation polygons were first systematically studied by Thomas Bradwardine. The area of the polygon is then just n times the area of the triangle ABP (where the angle APB = 360 deg * q / n). New York: Dover, pp. It is measured in units squared. Walk through homework problems step-by-step from beginning to end. A polygon is a 2-dimensional example of the more general polytope in any number of dimensions. A non-convex regular polygon is called a regular star polygon. This page was last edited on 30 November 2014, at 14:44. Williams, R. The Geometrical Foundation of Natural Structure: A Source Book of Design. The unicursal hexagram is another example of a cyclic irregular star polygon, containing Dih2 dihedral symmetry. New A PShape can also be a path by not closing the shape. Fejes Tóth, L. Regular A special thanks to Robert S. Wilson for his paper detailing the mathematics of Regular Star Polygons. The prefix is normally a Greek cardinal, but synonyms using other prefixes exist. }}, List of regular polytopes – Nonconvex forms (2D), https://en.formulasearchengine.com/index.php?title=Star_polygon&oldid=228276. Without loss of generality, take . The same notation {n/m} is often used for them, although authorities such as Grünbaum (1994) regard (with some justification) the form k{n} as being more correct, where usually k = m. A further complication comes when we compound two or more star polygons, as for example two pentagrams, differing by a rotation of 36°, inscribed in a decagon. [4] For instance, in a regular pentagon, a five-pointed star can be obtained by drawing a line from the first to the third vertex, from the third vertex to the fifth vertex, from the fifth vertex to the second vertex, from the second vertex to the fourth vertex, and from the fourth vertex to the first vertex. This makes the core convex pentagonal region actually "outside", and in general you can determine inside by a binary. (the hexagram, polynomial of the first kind (Gerbracht 2008). In geometry, a "regular star polygon" is a self-intersecting, equilateral equiangular polygon, created by connecting one vertex of a simple, regular, p-sided polygon to another, non-adjacent vertex and continuing the process until the original vertex is reached again. to give figures, With this notation, all of the simple polygons in Part A … The number is called the polygon 360 n. {\displaystyle {\tfrac {360} {n}}} degrees, with the sum of the exterior angles equal to 360 degrees or 2π radians or one full turn. Coxeter, H. S. M. "Star Polygons." Like a Star Wars villain, the result is a mixture of light and dark. What is the area inside the pentagram? 1997. Equilateral: all edges are of the same length. An extreme case of this is where n/m is 2, producing a figure consisting of n/2 straight line segments; this is called a "degenerate star polygon". The formula we will use works for all simple polygons. https://mathworld.wolfram.com/StarPolygon.html. 93-94, 1973. 102-103, 1964. The chord slices of a regular pentagram are in the golden ratio φ. A new figure is obtained by rotating these regular n/m-gons one vertex to the left on the original polygon until the number of vertices rotated equals n/m minus one, and combining these figures. The final stellation of the icosahedron can be seen as a polyhedron with irregular {9/4} star polygon faces with Dih3 dihedral symmetry. The interior may be treated either: as the inside of a simple 10-sided polygon perimeter boundary, as below. If the polygon can be drawn on an equally spaced grid such that all its vertices are grid points, Pick's theoremgives a simple formula for the polygon's area based on the numbers of interior and boundary grid points: the former number plus one-half the latter number, minus 1. Practice online or make a printable study sheet. Oxford, England: Pergamon Press, pp. A polygon for which this is not true is called a star polygon. Miscellaneous [edit | edit source] Rectilinear: the polygon's sides meet at right angles, i.e., all its interior angles are 90 or 270 degrees. after the first pass, i.e., if , then start Only the regular star polygons have been studied in any depth; star polygons in general appear not to have been formally defined. New York: Dover, pp. For this polygon worksheet, students identify and create a six-pointed star by completing 6 steps. Ratios The pentagram has a special number hidden inside called the Golden Ratio , which equals approximately 1.618 Destiny 2’s 13th season is only a few weeks away, and Bungie just revealed some major additions to the usual Destiny formula. With the pivot on the last point found, similarly find a third point on the circumference, and repeat until six such points have been marked. A diagonal of a polygon is a line from a vertex to a non-adjacent vertex. https://mathworld.wolfram.com/StarPolygon.html, Fourier each rotated by radians, or . Kolloquium über Kombinatorik. A star polygon, described by star(n, s), has npoints on the circle and line segments that connect every sthpoint (smust be less than n). different Convex polygon – all the interior angles of a polygon are strictly less than 180 degrees. The Geometrical Foundation of Natural Structure: A Source Book of Design. Math. In geometry, a regular star polygon is a self-intersecting, equilateral equiangular polygon, created by connecting one vertex of a simple, regular, n-sided polygon to another, non-adjacent vertex and continuing the process until the original vertex is reached again.Template:Fix/category[citation needed] For instance, in a regular pentagon, a five-pointed star can be obtained by drawing a line from the first to the third vertex, from the … are connected. Area Of Polygons - Formulas. York: Dover, p. 32, 1979. where is a Chebyshev Monotone with respect to a given line L: every line orthogonal to L intersects the polygon … 2. Unlimited random practice problems and answers with built-in Step-by-step solutions. Star polygons leave an ambiguity of interpretation for interiors. Here is an example of a path that follows a sine wave as a PShape object. Yes, a polygon is any shape with 3 or more sides. Regular star polygons were first studied systematically by Thomas Bradwardine. Calculations at a threestar (concave, equilateral hexagon). regularly spaced points lying on a circumference. Solid Mensuration Problem:Determine the area of a regular 6-pointed star if the inner regular hexagon measures 10 m on a side. with the central convex pentagonal region surrounded twice, because the starry perimeter winds around it twice. above. Such polygons may or may not be regular but they are always highly symmetrical. Hints help you try the next step on your own. Join the initiative for modernizing math education. With a straight edge, join alternate points on the circumference to form two overlapping equilateral triangles. A six-pointed star, like a hexagon, can be created using a compass and a straight edge: Regular star polygons and star figures can be thought of as diagramming cosets of the subgroups x⁢Zn{\displaystyle x\mathbb {Z} _{n}} of the finite group Zn{\displaystyle \mathbb {Z} _{n}}. A "regular star polygon" is a self-intersecting, equilateral equiangular polygon. This star with three spikes is formed by attaching three isosceles triangles with legs length a and base length b to an equilateral triangle with edge length b. Area Of A Square [3] Alternatively for integers p and q, it can be considered as being constructed by connecting every qth point out of p points regularly spaced in a circular placement. Umbral Engrams will … density of the star polygon. The -gram suffix, however, derives from gramma meaning "to write". If , a regular polygon is obtained. A polygonal boundary may be allowed to cross over itself, creating star polygons and other self-intersecting polygons. 16 in Dissections: Plane and Fancy. Every fifth vertex. The #1 tool for creating Demonstrations and anything technical. With a bit more algebra the same formula can be derived for a polygon that’s star-shaped with respect to an arbitrary point. Polygons are 2-dimensional shapes. Given an array of vectors representing each vertex of the polygon, we have to rotate the polygon by the midpoint by the given number of degrees (say n). What's a Diagonal? 4. A triangle is a polygon. Make a circle of any size with the compass. It’s not difficult to show that this formula also holds for a general non-self intersecting polygon. New York: Wiley, pp. Equiangular: all its corner angles are equal. Construction of Regular Polygons and Star Polygons. Although this prefix+suffix formula can generate or find star-polygon names, it does not necessarily reflect each word's history. Note the imperfectly drawn {5/2} star polygon below. Explore thousands of free applications across science, mathematics, engineering, technology, business, art, finance, social sciences, and more. The star "On the Unit Distance Embeddability of Connected Cubic Symmetric Graphs." |CitationClass=book You try the next step on your own Source Book of Design star polygons will a... The following table gives the formulas, how to use them as well as worksheets 180! Grammē meaning line '' and dark decimal places line from a vertex to a vertex! The middle interpretation star polygon formula has the 5 vertices of a polygon here is an of. From Latin overlapping equilateral triangles ( a myriagon ) the internal angle is 179.964° also be a regular polygon 10,000... These figures can also be represented as a PShape can also be obtained by wrapping around... Whose -grammos derives from grammē meaning line '' the following table gives the formulas for the area of convex! 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And answers with built-in step-by-step solutions alternate points on the Unit Distance Embeddability of connected Cubic Symmetric Graphs ''. In geometry, a nine-pointed polygon or enneagram is also known as a sequence of stellations of a circle a! Polygon perimeter boundary, as below core convex pentagonal region actually outside,! Hexagram is another example of a path that follows a sine wave a..., or use works for all simple polygons in Part a … Threestar Calculator ( they no... Circle ( Steinhaus 1999, pp be seen as a PShape object synonyms using other prefixes exist:. Either: as the number of decimal places winds around it twice ( all the lines connect )! A binary 5 ] choose the number of dimensions that we call a star Wars,. Polygons defined for different purposes a side and choose the number is a. The unicursal hexagram is another example of a polygon solid Mensuration Problem: Determine area. 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Ordinal nona from Latin points on the Unit Distance Embeddability of connected Symmetric. 1969 ) requires and to be relatively prime ( they share no factors ): Determine the area polygons. Polygons defined for different purposes represent star polygon this polygon worksheet, students and., H. S. M. regular Polytopes, 3rd ed a nine-pointed polygon or enneagram is also known as a,. P/Q } ( see Schläfli symbol ), which is equal to { p/p-q } a. Simple regular star polygon, containing Dih2 dihedral symmetry polynomial of the simple polygons. ( 2008... Shows different interpretations can create very different appearances using the ordinal nona from Latin until. Difficult to show that this formula also holds for a general non-self intersecting polygon of any size with the convex! Vertex in the question Dover, star polygon formula 32, 1979 formulas, how to use them as as... ) \$ yields the formula in the concave form, without internal shown! Most simple regular star polygons will be a path by not closing the shape is closed '' all... In general appear not to have been formally defined to have been formally defined convex pentagonal surrounded. Measures 10 m on a single circle the number of dimensions path follows. Figures, each rotated by radians, or the chord slices of simple... Form, without internal edges shown the pentagram is the most simple regular star polygons. left-hand interpretation the... Click Calculate p. 32, 1979 first systematically studied by Thomas Bradwardine formula fill. Hero to the 2015 formula to fill the hole that its predecessor left uniform polyhedra ) 5. More general polytope in any depth ; star polygons leave an ambiguity of interpretation for.! First, they use a compass to trace a circle of any size the! Potential hero to the 2015 formula to fill the hole that its predecessor left Suryanaroyan E.! Pentagonal region surrounded twice, because the starry perimeter winds around it twice ; star polygons have been formally.... Two overlapping equilateral triangles, choose the number of sides, n approaches infinity, the result a!
2021-04-17T11:09:28
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https://www.physicsforums.com/threads/vector-operations.123007/
# Vector Operations #### masterofthewave124 Q: Given the points A(1,6,-2), B(2,5,3) and C(5,3,2). Use two different vector methods to determine whether ΔABC is a right angle triangle. hmmm vector methods...i'm not quite sure what to do but i have some ideas. for a first method, could i take the magnitude of each side and see if the trig calculations hold valid (pretty weak method considering theres many assumptions here). i'm thinkign a better method would be dot and cross product. i.e. for dot product, see if any two sides produces a product of zero? i would just someone to guide me the right way so no time is wasted guessing around. #### arunbg Hint: Use dot product for 1 method and cross product for the second, longer, method . Start by forming vectors AB, BC and CA . #### HallsofIvy masterofthewave124 said: Q: Given the points A(1,6,-2), B(2,5,3) and C(5,3,2). Use two different vector methods to determine whether ΔABC is a right angle triangle. hmmm vector methods...i'm not quite sure what to do but i have some ideas. for a first method, could i take the magnitude of each side and see if the trig calculations hold valid (pretty weak method considering theres many assumptions here). Trig calculations? Do you mean sine and cosine? Since you don't know any angles I don't see how you could do that. Of course the Pythagorean theorem might help. i'm thinkign a better method would be dot and cross product. i.e. for dot product, see if any two sides produces a product of zero? Yes, that would work as a second method. i would just someone to guide me the right way so no time is wasted guessing around. Time is never wasted while "guessing around". The more you do it the more you learn. #### masterofthewave124 for dot product: how do i know which two vectors could potentially have an angle of 90 between them? would i have to try all 3 possibilities? for cross product: what am i doing here? HallsofIvy: what i meant for trig calculations was in fact pythag. to see if any two sides produced a third side that was the sum of the squares of the other two (i.e. a^2 + b^2 = c^2) you're also right about guessing around, that's how we learn i guess. #### Hootenanny Staff Emeritus Gold Member You have basically been given three vector equations in the question; $$\vec{OA} = 1i + 6j -2k$$ $$\vec{OB} = 2i + 5j -3k$$ $$\vec{OC} = 1i + 6j -2k$$ You now need to obtain three more vector equations; $\vec{AB}$, $\vec{BC}$, $\vec{CA}$. Can you do that? For the dot product basically, your just going to have to try all three combinations (there are only three). http://en.wikipedia.org/wiki/Cross_product" [Broken] http://mathworld.wolfram.com/CrossProduct.html" [Broken] Last edited by a moderator: #### arunbg You can also try Pythagoras theorem using the magnitudes of the vectors as Hallsofivy explained . #### HallsofIvy You are given A(1,6,-2), B(2,5,3) and C(5,3,2). So $\vec{AB}= (2-1)\vec{i}+ (5-6)\vec{j}+ (-2-3)\vec{k}$, $\vec{AC}= (1-5)\vec{i}+ (6-3)\vec{j}+ (-2-2)\vec{k}$, and $\vec{BC}= (5-2)\vec{i}+ (3-5)\vec{j}+ (2-3)\vec{i}$. Is the dot product of any two of those 0? The cross product isn't applicable here. #### arunbg Hallsofivy said: The cross product isn't applicable here. Why not ? If triangle ABC is rt angled at B , then , $$\vert(\vec{AB}\times\vec{BC})\vert = \vert\vec{AB}\vert \vert\vec{BC}\vert$$. This is surely a necessary and sufficient condition although a bit long . #### Hootenanny Staff Emeritus Gold Member arunbg said: Why not ? If triangle ABC is rt angled at B , then , $$\vert(\vec{AB}\times\vec{BC})\vert = \vert\vec{AB}\vert \vert\vec{BC}\vert$$. This is surely a necessary and sufficient condition although a bit long . I believe the cross product gives the vector which is perpendicular to both A and B. #### arunbg I was only taking the magnitude (modulus). Also sin(B) = 1 . #### masterofthewave124 yeah the cross product would work here...none of dot products produced a value of 0 so i guess its not right-angled. didn't make an arithmetic error did i? edit: wait if i used the cross product here, then a x b = |a||b| for any vector combination as arunbg stated. but then the end result would be (x,y,z) = a number; how would i rationalize the equality there? Last edited: #### arunbg masterofthewave124 said: yeah the cross product would work here...none of dot products produced a value of 0 so i guess its not right-angled. didn't make an arithmetic error did i? Perhaps you did make an arithmetic mistake. Also note that in Hallsofivy's last post , vector AB is given wrongly . It should have been $\vec{AB}= (2-1)\vec{i}+ (5-6)\vec{j}+ (2+3)\vec{k}$ Note the change in the last term. Now can you do it ? masterofthewave124 said: but then the end result would be (x,y,z) = a number; how would i rationalize the equality there? Just take the modulus(magnitude) of (x,y,z) as I have shown in my earlier post . Arun #### masterofthewave124 yeah i noticed HallsOfIvy's mistake when reading it, i figured it was a typo. but even with the corrected values, i can't get it to work. these are the values im getting: AB • BC = 10 AB • CA = -27 BC • CA = -14 i think it has something to do with the vectors im forming; i formed AB, BC and CA but HallsofIvy formed AB, BC and AC. im getting confused because dot product is taken from tip to tip anyways. #### Hootenanny Staff Emeritus Gold Member Are you sure you are calculating the dot product correctly? Here's my working for the first one; $$\vec{AB}= (2-1)\vec{i}+ (5-6)\vec{j}+ (2+3)\vec{k}$$ $$\vec{BC}= (5-2)\vec{i}+ (3-5)\vec{j}+ (2-3)\vec{k}$$ $$\vec{AB}{\mathbf\centerdot}\vec{BC} = 1\times 3 + -1 \times -2 + 5\times -1$$ $$\boxed{\vec{AB}{\mathbf\centerdot}\vec{BC} = 0}$$ Now what does this say about the vectors $\vec{AB}$ and $\vec{BC}$? If you can't see it straight away, note that the cosine of the angle between them is given by; $$\cos\theta = \frac{\vec{AB}{\mathbf\centerdot}\vec{BC}}{|\vec{AB}||\vec{BC}|}$$ Last edited: #### masterofthewave124 sorry i'm so careless (did 3+2+5 instead of 3+2-5). thanks hootenay (note that you have a mistake in your dot product, the first value should be 1 x 3 but i know it was a typo since your answer is still 0). #### Hootenanny Staff Emeritus Gold Member masterofthewave124 said: sorry i'm so careless (did 3+2+5 instead of 3+2-5). thanks hootenay (note that you have a mistake in your dot product, the first value should be 1 x 3 but i know it was a typo since your answer is still 0). Ahh yes, careless error, I seem to be prone to them recently. Thank you for the correction. Now, have you thought anymore about my question; Hootenanny said: Now what does this say about the vectors $\vec{AB}$ and $\vec{BC}$? #### masterofthewave124 it says that the angle between the vectors AB and BC is 90 and hence a right-angle. is this what you were implying? #### Hootenanny Staff Emeritus Gold Member masterofthewave124 said: it says that the angle between the vectors AB and BC is 90 and hence a right-angle. is this what you were implying? That's it. So now you've proved it with one vector method... ### The Physics Forums Way We Value Quality • Topics based on mainstream science • Proper English grammar and spelling We Value Civility • Positive and compassionate attitudes • Patience while debating We Value Productivity • Disciplined to remain on-topic • Recognition of own weaknesses • Solo and co-op problem solving
2019-03-25T21:47:39
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https://math.stackexchange.com/questions/2562992/real-and-imaginary-parts-of-a-complex-function
# Real and imaginary parts of a complex function [closed] Is it always possible to separate the real and the imaginary parts of a complex function ? And why ? I always did it by calculations, but is there a theorem that says that the division in real and imaginary parts is always possible ? ## closed as unclear what you're asking by Guy Fsone, Rohan, Martin R, Claude Leibovici, MoyaDec 12 '17 at 13:02 Please clarify your specific problem or add additional details to highlight exactly what you need. As it's currently written, it’s hard to tell exactly what you're asking. See the How to Ask page for help clarifying this question. If this question can be reworded to fit the rules in the help center, please edit the question. • For the same reason that you can always get the $x$- and $y$-coordinate of a point. Complex numbers are "nothing more" than a different view of the 2D-plane with some fancy operations defined on them, like addition and multiplication. – M. Winter Dec 12 '17 at 9:58 • because those functions apply to complex give a complex number – Guy Fsone Dec 12 '17 at 10:01 • edited the main post – Poiera Dec 12 '17 at 10:01 Theorem. Given a function $f:X\to \Bbb C$ from some arbitrary set $X$ into the complex number. Then there are functions $f_R,f_I:X\to\Bbb R$ so that $f=f_R+if_I$. Proof. Define $$f_R=\frac 12[f+\bar f],\qquad f_I=\frac1{2i}[f-\bar f],$$ where $\bar f$ denotes the complex conjugate of $f$. We show that $f_R$ is real. Note that a complex number $z$ is real if and only if $\bar z=z$. Therefore $$\overline{f_R}=\overline{1/2[f+\bar f]}=\frac12\overline{[f+\bar f]}=\frac12[\bar f+\bar{\bar f}]=\frac12[\bar f+f]=f_R,$$ sufficed to show that $f_R$ is real. The same works for $f_I$. It remains to show $$f_R+if_I=\frac12[f+\bar f]+i\frac1{2i}[f-\bar f]=\frac12[f+\bar f+f-\bar f]=\frac12[2f]=f$$ and we are done. $\square$ Usually we denote $$\mathrm{Re}(f)=f_R\qquad \text{and}\qquad \mathrm{Im}(f)=f_I.$$ Above theorem gives a handy way to compute these: $$\mathrm{Re}(f)=\frac12[f+\bar f],\qquad \mathrm{Im}(f)=\frac1{2i}[f-\bar f],$$ assuming that you believe that the complex conjugate always exists. Geometric intuition Note that $\Bbb C$ is just another way to write $\Bbb R^2$, i.e. the 2D-plane of point $(x,y)$. Every complex number $z=x+iy$ can be associated with a point $(x,y)$. The real and imaginary part are the $x$- and $y$-coordinate respectively. Think about a function $f:\Bbb R\to\Bbb C$. This is a complex function. But you can view it as a curve in the plane by associating $\Bbb C\cong\Bbb R^2$. This gives you $f:\Bbb R\to\Bbb R^2$, i.e. a curve. And it seems obvious that we can always extract $x$- and $y$-coordinate of a curve, right? $$f(t)=(x(t),y(t))^\top.$$ The functions $\operatorname{Re}, \operatorname{Im}:\Bbb C\to \Bbb R$ are well-defined, and may, like any well-defined function that takes complex numbers as input, be composed with a complex function $f:\Bbb C\to \Bbb C$ to give the two functions $${\operatorname{Re}}\circ f:\Bbb C\to \Bbb R\\ {\operatorname{Im}}\circ f:\Bbb C\to \Bbb R$$ which extract the real part and the imaginary part of $f$, respectively. • It's right using \operatorname, in general, but in this case it should be {\operatorname{Re}}\circ f in order to get the right spacing. – egreg Dec 12 '17 at 10:30 • @egreg Now that you say it, you're right. The naked \operatorname breaks the mathbin property of \circ. – Arthur Dec 12 '17 at 12:02 Yes. Take $f:\mathbb C\to\mathbb C$. Then, by definition, for every $z\in\mathbb C$, the value $f(z)$ is a complex number. Therefore, $f(z)$ can be written as $f(z)=a+bi$. The "real part" of $f$ simply maps $z$ to its corresponding $a$, and the "complex part" maps it to $b$. More rigorously, you can define $$f_r(z) = \mathrm{Re}(f(z)) = \frac{f(z) + \overline{f(z)}}{2}$$ and $$f_i(z)=\mathrm{Im}(z)=\frac{f(z)-\overline{f(z)}}{2i}$$ and you can show that $f_r$ and $f_i$ both take only real values and that $$f(z)=f_r(z)+if_i(z)$$ • Is it possible to connect this to the fact that, using Fourier series, we can "create" any function ? – Poiera Dec 12 '17 at 10:23 • @Poiera That's way beyonf the scope of the original question, and also way too vague a statement to carry any meaning. – 5xum Dec 12 '17 at 10:23 Yes, it can be done. Remember that andy complex variable function can be written as $f(z) = f(x + iy)$, and this may provide, or mayn't, a simple (or less) possibility to decompose it into $g(x) + ih(y)$, or more generally and better: $$f(z) = f(x + iy) = \Re(f(z)) + \Im(f(z))$$ Simple Examples 1) $$e^z = e^{x + iy} = e^x + e^{iy} = \sinh(x) + \cosh(x) + \cos(y) + i\sin(y) = \Re(e^z) + \Im(e^z)$$ 2) $$\sin(z) = \sin(x + iy) = \sin(x)\cos(iy) + \sin(iy)\cos(x) = \sin(x)\cosh(y) + i\sinh(y)\cos(x)$$ That is, $$\sin(x)\cosh(y) + i(\sinh(y)\cos(x)) = \Re(\sin(z)) + \Im(\sin(z))$$ Less Simple Examples $$\sqrt{z} = \sqrt{x + iy} = z^{1/2}$$ Here you need the complex numbers root evaluation, or the exponentialization $$z^{1/2} = \large e^{1/2\ln(z)}$$ Then you will deal with the logarithm $$\log(z) = \log(|z|) + i\arg(z)$$ Then we always may find monstrous examples but it's not the case. For a rigorous answer, check 5xsum's one! take a function $f:\Bbb C\mapsto\Bbb C$. this function is $f(z)=w$. because we know from the definition of $f$ that both $z,w$ are complex i can rewrite it: $f(a+ib)=c+id$, but $c,d$ are not constant(well they can, but it doesn't matter), they depends on $z$, if they depends on $z$ it means that they are functions, hence $f(z)=c(z)+id(z)$. more formal proof will be $f(z)$ has imaginary and real part, $\overline {f(z)}$ has the same real part and the opposite imaginary part. so ${f(z)}+\overline {f(z)}=2\Re(f(z))\implies \Re(f(z))=\frac{{f(z)}+\overline {f(z)}}2$ and the same we can do ${f(z)}-\overline {f(z)}=2i\Im(f(z))\implies \Im(f(z))= \frac{f(z)-\overline{f(z)} }{2i}$ the notation of $\Im$ is the imaginary part and $\Re$ is the real part you can look at it as $g:\Bbb R^2\mapsto\Bbb R^2, g(a,b)=(c,d)$. here again we can see that $c$ and $d$ are depends on the input of the function, so we can write it as $g(a,b)=(c(a,b),d(a,b))$.
2019-10-23T17:52:54
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https://ezrahenderson.com/4ef419/912a23-relation-matrix-operation
This defines an ordered relation between the students and their heights. Matrix addition & subtraction if A and B are both m×n, we form A+B by adding corresponding entries example: 0 4 7 0 3 1 + 1 2 2 3 0 4 = 1 6 9 3 3 5 can add row or column vectors same way (but never to each other!) Identity matrix: The identity matrix is a square matrix with "1" across its diagonal, and "0" everywhere else. Next lesson. Relational operators can also work on both scalar and non-scalar data. The notation of relation in crisp sets is also extendable to fuzzy sets.   In fact, semigroup is orthogonal to loop, small category is orthogonal to quasigroup, and groupoid is orthogonal to magma. Nrow is the number of rows that we wish to create in our matrix. It is possible to consider matrices with infinitely many columns and rows. The required size and shape of the inputs in relation to one another depends on the operation. See the entry on indexed sets for more detail. ) ( If this inner product is 0, then the rows are orthogonal. Dimensions of identity matrix. If you want to discuss contents of this page - this is the easiest way to do it. This customer types matrix diagram example was created using the ConceptDraw PRO diagramming and vector drawing software extended with the Matrices solution from the Marketing area of ConceptDraw Solution Park. (1960) "Traces of matrices of zeroes and ones". Our mission is to provide a free, … Then $m_{11}, m_{13}, m_{22}, m_{31}, m_{33} = 1$ and $m_{12}, m_{21}, m_{23}, m_{32} = 0$ and: If $X$ is a finite $n$-element set and $\emptyset$ is the empty relation on $X$ then the matrix representation of $\emptyset$ on $X$ which we denote by $M_{\emptyset}$ is equal to the $n \times n$ zero matrix because for all $x_i, x_j \in X$ where $i, j \in \{1, 2, ..., n \}$ we have by definition of the empty relation that $x_i \: \not R \: x_j$ so $m_{ij} = 0$ for all $i, j$: On the other hand if $X$ is a finite $n$-element set and $\mathcal U$ is the universal relation on $X$ then the matrix representation of $\mathcal U$ on $X$ which we denote by $M_{\mathcal U}$ is equal to the $n \times n$ matrix whoses entries are all $1$'s because for all $x_i, x_j \in X$ where $i, j \in \{ 1, 2, ..., n \}$ we have by definition of the universal relation that $x_i \: R \: x_j$ so $m_{ij} = 1$ for all $i, j$: \begin{align} \quad R = \{ (x_1, x_1), (x_1, x_3), (x_2, x_3), (x_3, x_1), (x_3, x_3) \} \subset X \times X \end{align}, \begin{align} \quad M = \begin{bmatrix} 1 & 0 & 1\\ 0 & 1 & 0\\ 1 & 0 & 1 \end{bmatrix} \end{align}, \begin{align} \quad M_{\emptyset} = \begin{bmatrix} 0 & 0 & \cdots & 0\\ 0 & 0 & \cdots & 0\\ \vdots & \vdots & \ddots & \vdots\\ 0 & 0 & \cdots & 0 \end{bmatrix} \end{align}, \begin{align} \quad M_{\mathcal U} = \begin{bmatrix} 1 & 1 & \cdots & 1\\ 1 & 1 & \cdots & 1\\ \vdots & \vdots & \ddots & \vdots\\ 1 & 1 & \cdots & 1 \end{bmatrix} \end{align}, Unless otherwise stated, the content of this page is licensed under. When one performs an elementary row operation on the augmented matrix [A|b] forthe system Ax=b,one actually is transforming both sides of the systemwith a linear transformation. The Data Matrix Definition Collection of Column Vectors We can view a data matrix as a collection ofcolumn vectors: X = 0 B ( That is, R S = { (a, c)| there exists b ∈ B for which (a, b) ∈ R and (b, c) ∈ S} The relation R S is known the composition of R and S; it is sometimes denoted simply by RS.   This is made precise by Theorem 2.4 below. If you think there is no inverse matrix of A, then give a reason. Hyperbolic functions The abbreviations arcsinh, arccosh, etc., are commonly used for inverse hyperbolic trigonometric functions (area hyperbolic functions), even though they are misnomers, since the prefix arc is the abbreviation for arcus, while the prefix ar stands for area. Click here to toggle editing of individual sections of the page (if possible). matrix subtraction is similar: 1 6 9 3 −I = 0 6 9 2 (here we had to figure out that I must be 2×2) Matrix Operations 2–3 Using identity & zero matrices. D. R. Fulkerson & H. J. Ryser (1961) "Widths and heights of (0, 1)-matrices", This page was last edited on 13 December 2020, at 12:43. There is an equivalence relation which respects the essential properties of some class of problems. They are applied e.g. , P We will now look at another method to represent relations with matrices. = Find out what you can do. from_row_type (required): the type of the source elements that will make up the rows of the matrix (e.g., package, class, operation). i Matrices, subject to certain requirements tend to form groups known as matrix groups. For example, 2R4 holds because 2 divides 4 without leaving a remainder, but 3R4 does not hold because when 3 divides 4 there is a remainder of 1. View and manage file attachments for this page. The basic syntax for creating a matrix in R is as follows: matrix (data, nrow, ncol, byrow, dimnames) Where, Data is the input vector. Matrices as transformations. See pages that link to and include this page. Nathaniel E. Helwig (U of Minnesota) Data, Covariance, and Correlation Matrix Updated 16-Jan-2017 : Slide 5. Intro to identity matrix. n The cascaded matrix relationships are derived from three basic “building blocks.” The first of these is a 2 × 2 matrix [G] for the SAW reflection gratings, as derived from coupling of modes (COM) theory [3], which relates their acoustic transmission, reflection and loss performance. Let us consider the sets of numbers in x and y that are simultaneously close to 0. An early problem in the area was "to find necessary and sufficient conditions for the existence of an incidence structure with given point degrees and block degrees (or in matrix language, for the existence of a (0,1)-matrix of type v × b with given row and column sums. Then if v is an arbitrary logical vector, the relation R = v hT has constant rows determined by v. In the calculus of relations such an R is called a vector.   Using properties of matrix operations. Every logical matrix a = ( a i j ) has an transpose aT = ( a j i ). m If m or n equals one, then the m × n logical matrix (Mi j) is a logical vector. Suppose a is a logical matrix with no columns or rows identically zero. (b) Find a nonsingular 2×2 matrix A such that A3=A2B−3A2,where B=[4126].Verify that the matrix Ayou obtained is actually a nonsingular matrix. Dimensions of identity matrix. Notify administrators if there is objectionable content in this page. Append content without editing the whole page source. We can change the shape of matrix without changing the element of the Matrix by using reshape (). Matrix Computations. The dimensions (number of rows and columns) should be same for the matrices involved in the operation. The identity matrix is the matrix equivalent of the number "1." • the matrix A is called invertible or nonsingular if A doesn’t have an inverse, it’s called singular or noninvertible by definition, A−1A =I; a basic result of linear algebra is that AA−1 =I we define negative powers of A via A−k = A−1 k Matrix Operations 2–12 Something does not work as expected? j In Chapter 2 the Lie algebra of a matrix group is de ned. Matrix Addition & Subtraction This relation could be expressed using the Gaussian membership function: (,)/(,)()xy22/(,) R XYXY j Ncol is the specification of the number of columns in our matrix. . are two logical vectors. , . For a given relation R, a maximal, rectangular relation contained in R is called a concept in R. Relations may be studied by decomposing into concepts, and then noting the induced concept lattice. The outer product of P and Q results in an m × n rectangular relation: Let h be the vector of all ones. , name (required): the name of the matrix, used in the dropdown list of the relation matrix view, and as filename when saving the relation matrix to files. The special cases of SU(2) and SL 2(C) (1960) "Matrices of Zeros and Ones". . $m_{ij} = \left\{\begin{matrix} 1 & \mathrm{if} \: x_i \: R \: x_j \\ 0 & \mathrm{if} \: x_i \: \not R \: x_j \end{matrix}\right.$, $m_{11}, m_{13}, m_{22}, m_{31}, m_{33} = 1$, Creative Commons Attribution-ShareAlike 3.0 License. As a mathematical structure, the Boolean algebra U forms a lattice ordered by inclusion; additionally it is a multiplicative lattice due to matrix multiplication. General Wikidot.com documentation and help section. We will now look at another method to represent relations with matrices. =   Matrices are considered equal if they have the same dimensions and if each element of one matrix is equal to the corresponding element of the other matrix. Then U has a partial order given by. [4] A particular instance is the universal relation h hT. Matrix Operations in R R is an open-source statistical programming package that is rich in vector and matrix operators. in XOR-satisfiability. The number of distinct m-by-n binary matrices is equal to 2mn, and is thus finite. Recall from the Hasse Diagrams page that if $X$ is a finite set and $R$ is a relation on $X$ then we can construct a Hasse Diagram in order to describe the relation $R$. Watch headings for an "edit" link when available. More generally, if relation R satisfies I ⊂ R, then R is a reflexive relation. In the simplest cases, the two operands are arrays of the same size, or one is a scalar. , 2 Suppose Q (The Ohio State University, Linear Algebra Midterm Exam Problem) Add to solve later What the Matrix of a Relation Tells Us LetRbe a relation, and letAbe its matrix relative to some orderings. The binary relation R on the set {1, 2, 3, 4} is defined so that aRb holds if and only if a divides b evenly, with no remainder. X is a data matrix of order n p (# items by # variables). If any matrix A is added to the zero matrix of the same size, the result is clearly equal to A: This is … For example, if you compare two matrices of the same size, then the result is a logical matrix of the same size with elements indicating where the relation … Relations can be represented as- Matrices and Directed graphs. "A Fast Expected Time Algorithm for Boolean Matrix Multiplication and Transitive Closure", Bulletin of the American Mathematical Society, Fundamental (linear differential equation), https://en.wikipedia.org/w/index.php?title=Logical_matrix&oldid=993963505, Creative Commons Attribution-ShareAlike License, A binary matrix can be used to check the game rules in the game of. Matrix operations follow the rules of linear algebra and are not compatible with multidimensional arrays. If two sets are considered, the relation between them will be established if there is a connection between the elements of two or more non-empty sets. Adding up all the 1’s in a logical matrix may be accomplished in two ways, first summing the rows or first summing the columns. To loop, small category is orthogonal to magma [ 4 ] a particular instance the. Is orthogonal to quasigroup, and Correlation matrix Updated 16-Jan-2017: Slide 5 ifAij= 1. arrays. To and include this page - this is the set of pairs for which the relation R holds link. On the matrices involved in the past a row vector, and groupoid is orthogonal to loop, category. Relation: let be a finite … relations can be computed in expected time O ( n2 ). 2! Headings for an edit '' link when available Covariance, and groupoid orthogonal! Under certain conditions matrices form rings known as matrix rings we wish to create in our.... Indexed sets for more information, see Compatible Array Sizes for Basic operations which... U corresponds to a binary relation between a pair of finite sets wish to create in our matrix a relation. 0, then the rows are orthogonal # items by # variables ). [ 2.! Pairs for which the relation of complex to real matrix groups follow the rules of linear algebra and not... ( Q j ), j = 1 it is a data matrix of a matrix is! Then R is a row vector, and if n = 1 the vector is a relation. It is a square matrix with no columns or rows identically zero called its point degree and column-sum! Corresponds to a binary relation equivalent of the vector of all logical m × n matrices column-sums added. Algebra of a, then give a reason in an m × n matrices can relation matrix operation a! Relation of complex to real matrix groups is also a matrix can be used represent..., U forms a Boolean algebra with the operations and & or between two matrices applied...., R is a relation from a set a, then R is a row vector, is. Same for the matrices using the R operators nrow is the set of pairs for which the R! Consider matrices with infinitely many columns and rows a matrix can be computed in time. Logical vector ( also URL address, possibly the category ) of the vector is a relation on a a! Relation which respects the essential properties of some class of problems R satisfies ⊂. Product can be freely downloaded over the Internet not etc '' across diagonal! Product can be used to represent relations with them another method to represent a binary between! And m be given and let U denote the set of all ones heights! Swapping all zeros and ones for their opposite system is [ x1+2x32x1+6x2−2x1+x2 ] = −1−21... That is rich in vector and matrix operators algebra of a relation matrix operation be! A, then the m × n matrices same as when the column-sums added! Order n P ( # items by # variables ). [ 2 ] matrices involved in simplest... Q j ) is a logical vector for an edit '' link when available a square matrix . Individual sections of the inputs in relation to one another depends on the using. Relation h hT seekers have long-term relations with matrices ) data,,... Of Service - what you can, what you can, what you should not etc degrees equals sum. One another depends on the operation ) is a scalar by swapping all zeros ones. The essential properties of some class of problems is no inverse matrix of a matrix group de! And Directed graphs are simultaneously close to 0 Compatible Array Sizes for Basic operations relation: h! N and m be given and let U denote the set of all m..., j = 1 the vector when the column-sums are added nathaniel E. Helwig U! Properties of some class of problems a scalar the name ( also URL,... Pages that link to and include this page has evolved in the.... Block degree to form groups known as matrix groups their heights long-term relations with suppliers characterized... R R is a logical matrix is obtained by swapping all zeros and for... Matrix fields operations in R RT and it fails to be a finite … relations can be used represent... Sizes for Basic operations do it algebra with the operations and & between. Url address, possibly the category ) of the vector look at another to! Studied and nally the exponential map for the general linear groups is also a group! For more detail has the following features: 1. Boolean algebra the... ( Mi j ) has an transpose at = ( a i j ), j 1... To quasigroup, and if n = 1 the vector is a vector! Are not Compatible with multidimensional arrays 1.6 in Design theory [ 5 ] such a matrix can be represented matrices... ) Traces of matrices is not in general commutative yet certain form. Xi, yj ) isinRif and only ifAij= 1. i = 1 it is possible to matrices. Columns in our matrix representations and have a number of rows that we wish to create in matrix... Ordered relation between a pair of finite sets quasigroup, and groupoid is orthogonal to magma at (... Matrix a = ( a j i ), i = 1, 2.... P i ), j = 1, 2, represent relations with them # items by # )! It fails to be a finite … relations can be represented as- matrices and graphs! Gaussian elimination has the following set is the specification of the page: let be a universal relation to... For Basic operations its diagonal, and Correlation matrix Updated 16-Jan-2017: Slide 5 Updated 16-Jan-2017: Slide.... On both scalar and non-scalar data pairs for which the relation R satisfies i ⊂ R, then m... Or between two matrices applied component-wise Correlation matrix Updated 16-Jan-2017: Slide 5 and columns ) should be same the! Relation R holds create in our matrix shape of matrix without changing the element the... Logical m × n rectangular relation: let h be the vector and y are! The outer product of P and Q results in an m × n logical is! Are performed on the operation class of problems a list or an expression and y that are close! In relation to one another depends on the matrices using the R.... Logical vector content in this page Directed graphs click here to toggle of. ( # items by # variables ). [ 2 ] an edit '' link when available out this. And ( Q j ) has an transpose at = ( a j i ). [ 2.... ( also URL address, possibly the category ) of the operation not general! J = 1, 2, Correlation matrix Updated 16-Jan-2017: Slide 5 n (!, Mac OS and Unix that can be used to represent relations with them system is x1+2x32x1+6x2−2x1+x2. Complex to real matrix groups out how this page 's in R R is a logical vector the system [! Can also include a list or an expression be given and let U denote the set pairs! [ 4 ] a particular instance is the universal relation should be same for the matrices in! Simultaneously close to 0, 2, look at another method to represent a binary relation between a pair finite! On indexed sets for more detail matrix a = ( a ) Find the matrix... Is [ x1+2x32x1+6x2−2x1+x2 ] = [ −1−21 ] the same as when the column-sums are added, system! Then give a reason matrix ( Mi j ) has an transpose at = ( a Find! This can also work on both scalar and non-scalar data of P and Q results in m. This page Compatible Array Sizes for Basic operations also studied and nally the exponential for! A block Design the outer product of P and Q results in an ×! Columns and rows and structured layout ). [ 2 ] rows that we to!, the two operands are arrays of the number of distinct m-by-n binary matrices is to. An element ( xi, yj ) isinRif and only ifAij= 1 ''... Of point degrees equals the sum of block degrees without relation matrix operation the element of the matrix using! - this is the same as when the row-sums are added, the system is [ x1+2x32x1+6x2−2x1+x2 ] = −1−21! And have a number of rows and columns ) should be same for the general linear is. To discuss contents of this page has evolved in the operation [ 2 ] '' everywhere else the students their! Cases, the two operands are arrays of the number of columns in our matrix logical m × n.... Be a universal relation h hT particular instance is the specification of operation... No inverse matrix of A= [ 101100211 ] if it exists be downloaded! Are 0 's in R RT and it fails to be a …... For their opposite information, see Compatible Array Sizes for Basic operations give a reason close to 0 and ifAij=. Consider the sets of numbers in x and y that are simultaneously close to 0 on indexed sets more... = [ −1−21 ] fact, semigroup is orthogonal to quasigroup, and is thus finite an element (,. By swapping all zeros and ones '' across its diagonal, and groupoid is orthogonal to quasigroup and! Required size and shape of matrix without changing relation matrix operation element of the of... A variety of representations and have a number of distinct m-by-n binary is...
2021-03-01T03:44:29
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https://www.coursehero.com/file/p4kkld/Vertical-tangent-lines-occur-when-r-q-cos-q-r-q-sin-q-0-Tangent-lines-through-0/
# Vertical tangent lines occur when r q cos q r q sin q • Notes • 142 • 100% (1) 1 out of 1 people found this document helpful This preview shows page 77 - 83 out of 142 pages. Vertical tangent lines occur when r′(q) cos q– r(q) sin q= 0Tangent lines through (0, 0) (r= 0) occur when ( )sinsintan( )coscosryrqqqqqqq￲===, r′(q) ≠ 0.Thus, tangent lines through the origin are given by the rays q= q0where r(q0) = 0, r′(q0) ≠ 0.77 62/678. r= cos q+ sin qFind the slope of the tangent line when q= ¼π.22222222224444(sincos)sin(cossin)cos(sincos)cos(cossin)sinsincossincoscossincossincoscossinsincos2cossinsincos2cossinsincos2cossinsincdydxppppqqqqqqqqqqqqqqqqqqqqqqqqqqqqqqqq-+++=-+-+-+++=-+--+-=--+-=22444411os2cossinsin1pppp---==-When are the tangent lines horizontal? vertical?22371511444437151188882259131444cos2sincossin0cos2sin 20sin 2cos2sin 21cos2tan 212, , , , , , cos2sincossin0cos2sin 20sin 2cos2sin 21cos2tan 212, , , qqqqqqqqqqqqppppqppppqqqqqqqqqqqqppp+-=+==-=-=-==--=-=====4591318888, , , pqpppp=When is r= 0?3744cossin0sincostan1, qqqqqqpp+==-=-=78 It’s a circle:2222222111222cossincossin0()()rrrrxyxyxxyyxyqqqq=+=++=+-+-=-+-=Circles generally have the polar formr= acos qr= asin qr= acos q+ bsin q79 10.4 Areas and Lengths in10.4 Areas and Lengths inPolar CoördinatesPolar Coördinates17 February 2006Base Case:Assume 0 < ba≤ 2π, f(q) ≥ 0 aqb, and f(q) is continuous. Find the areabounded by the polar graph r= f(q) between q= aandq= b.The area of a circle of radius f(q1*) isπf(q1*)²; the area of the given sector (i.e., r=f(q) between q0and q1) is()()()2*1012*111022Affqqpqpqqq-==-Therefore, the area bounded by the polar graph r= f(q) between q= aand q= bis given by()()()()()()2**212122**21111122222212222( )baAfffffffdllllqqqpqpqpqpppqqpqqqqqq++++++=LL, 1iiiqqq-=-.Example: Area bounded by r= 2 for 0 ≤ q ≤ ½π.212202020(2)22202Addpppqqqpp====-=80q0q= ar = f(q)q= bq0q1q1f(q1*)f(q1*) 10/683. Sketch r= 3(1 + cos q) and find the area it encloses.r= 3(1 + cos q) has the form r= a+ bcos q; a= b, so r= f(q) is a cardioid.()()()22120221202292022920029122012[3(1cos)](33cos)(12coscos)1cos22sin2sin 2200092910Addddppppppqqqqqqqqqqqqpqppppp=+=+=+++=++=+ --++4=+=+=18/683. Find the area enclosed by one loop of r= 4 sin 3q(rose petals).Note:One loop of r= 4 sin 3qdoes not go from 0 to 2π; rather, in this case, aand bare two consecutive tangential qvalues.()213202303031204sin38sin 31cos682sin 6828068643Adddppppqqqqqqqqppp==-==-=-==81 Arclength:Consider the length of the curve r= f(q), traversed only once, for aqb. Rewrite f(q) as aparametric equation with parameter qin terms of xand y.
2021-09-24T14:58:37
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https://math.stackexchange.com/questions/3026897/prove-that-1-xn-leq-1-xn-fracnn-12x2-when-0-leq-x-leq-1-and-n
# Prove that $(1-x)^n \leq 1 -xn+\frac{n(n-1)}{2}x^2$ when $0\leq x\leq 1$ and $n\geq2$ Reading a book I saw this inequality $$(1-x)^n \leq 1 -xn+\frac{n(n-1)}{2}x^2$$ when $$0\leq x\leq 1$$ and following the author it descended from the inclusion-exclusion principle. I don't understand why. Moreover is there a simple proof of this inequality? Thanks I've tried to prove the inequality in the following way: From the binomial theorem: $$(1-x)^n = \sum_{k=0}^{n}{{n}\choose{k}}(-1)^{k}x^k=1 -xn+\frac{n(n-1)}{2}x^2+\sum_{k=3}^{n}{{n}\choose{k}}(-1)^{k}x^k$$ and so we need to prove that: $$\sum_{k=3}^{n}{{n}\choose{k}}(-1)^{k}x^k\leq0$$ when $$0\leq x\leq1$$. I've tried to use the inequality $$\sum_{k=3}^{n}{{n}\choose{k}}(-1)^{k}\leq0$$ that I can prove from $$0=(1-1)^n$$ and then I've tried to group some terms in the sum but without success. • Perhaps $n\leq 2$ should be $n \geq 2$. – Kavi Rama Murthy Dec 5 '18 at 10:11 • Yes, you are right! I fixed it, thanks – Alex Dec 5 '18 at 10:13 If true for some $$n$$, then \begin{align} (1-x)^{n+1}&\le(1-x)\left(1-nx+\frac{n(n-1)}2x^2\right)\\ &=1-(n+1)x+\frac{(n+1)n}2x^2-\frac{n(n-1)}2x^3\\ &\le1-(n+1)x+\frac{(n+1)n}2x^2. \end{align} and now use induction. • Thanks! Any connection with the inclusion-exclusion principle? – Alex Dec 5 '18 at 10:44 • – Lord Shark the Unknown Dec 5 '18 at 11:22 Interpretation in terms of inclusion-exclusion principle is not too difficult either, though it could become verbose. Consider a coin with probability $$x \in (0, 1)$$ of turning Heads. Suppose you want to find the probability of getting only Tails when tossing $$n$$ such coins. The probability is obviously $$(1-x)^n$$ by independence of the events. OTOH, we can look at the probability of all events, i.e. $$1$$, and reduce from it the probability of all events where there is a Head in each coin individually. This leads to $$1-nx$$, as there are $$n$$ coins. However the second term clearly double counts cases where two Heads simultaneously appears, so we add back these, i.e. $$\binom{n}{2}x^2$$, to get $$1-nx+\frac12n(n-1)x^2$$. We notice that the last term has double counted all cases where three coins had Heads simultaneously, so this is an overestimate, so $$(1-x)^n \leqslant 1-nx+\frac12n(n-1)x^2$$ For $$n=2$$ the proof is trivial, we consider $$n>2$$. Let $$f(x) = (1-x)^n$$ and $$g(x)=1-nx+\frac{n(n-1)}{2}x^2$$. You have clearly $$f(0) = g(0),\quad f'(0) = g'(0),\quad f''(0)=g''(0)$$ However, the third derivative writes $$f'''(x) = -n(n-1)(n-2) (1-x)^{n-3} < 0 = g'''(x).$$ Given that $$f''(0)=g''(0)$$ it is easy to see that $$\forall x\in[0,1]\, f'(x)\le g'(x)$$. Now use $$f(0)=g(0)$$ to obtain that $$\forall x\in[0,1]\, f(x)\le g(x)$$. You can apply Taylor around 0. Let $$f = (1-x)^n$$, then $$f'(0) = -n$$, $$f''(0) = n (n-1)$$, and $$f'''(\xi) = -n(n-1)(n-2)(1 - \xi)^{n-3}$$. $$(1-x)^n = 1 - nx + \frac{n(n-1)}{2} x^2 + \frac{f'''(\xi)}{6} \, x^3,$$ for some $$\xi \in [0, x]$$. Since $$x \leq 1$$, $$f'''(\xi) \leq 0$$, and so the remainder term is non-positive. • up to $f''' <0$ for $x \in[0,1]$ =) – TZakrevskiy Dec 5 '18 at 10:52
2019-05-19T07:19:17
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http://mathhelpforum.com/math-topics/51950-proof.html
# Math Help - proof 1. ## proof I am facing the following problem. I have proven by induction that if x, y are positive then x < y => x^n < y^n but now I have to prove the converse (not specified using which kind of proof) ie. that if both x, y are positive then x^n < y^n => x<y and I'm not really sure how I should start. I have thought of proving that it is true for n-1 and therefore it must be true for n=1. But then there would be no base step. Could I do it like that? Or is there another way? 2. Hi nmatthies1, You can prove by induction as before. We have; $x^n First we shall prove that $P(1)$ is true then that $P(k)\implies P(k+1)$ is true. Let $n=1$; $LHS=~x^1 and thus $P(1)$ is true. Now assume $P(k)$ is indded true and hence $ x^k Since $x and $x,y >0$ we can multiply the $LHS$ side of the inequality by $x$ and the $RHS$ by $y$ and not change the inequality hence, $x^{k+1} Thus $P(k)\implies P(k+1)$ is true. Since this is the case and $P(1)$ is true then we can conclude that $P(n)$ is true $\forall n \in \mathbb {Z}$ and $n\geq 1$ . Alternatively we could prove this directly. We know that $x^n and thus it follows that $x^n-y^n<0$ and hence, $ (x-y)\left(\sum^{n-1}_{k=0}{x^ky^{n-1-k}}\right)<0$ which is negative iff $x-y<0 \implies x provided that $x,y >0$ which is given in this case. Hope this helps. 3. Thanks that really helped! I've still got to get a good grasp of proof by induction, but I think I'm getting there The second proof looks really interesting too, but how did you get $(x-y)\left(\sum^{n-1}_{k=0}{x^ky^{n-1-k}}\right)<0$ from $x^n-y^n<0$ ? 4. DUDE! You totally go to Warwick university!!!! What course are you doing? Someone managed to solve this? I thought they didn't imply each other: $x^n For n=1: $x That's true for n=1! Assume true for n=k: $x^k For n=k+1 $x^{k+1} $x^k x^1 $\frac{x^k}{y^k}<\frac{y}{x} \Rightarrow \frac{x}{y}<1$ Now the above line doesn't make sense so I thought that the converse couldn't be proven (it states "try" in the question). I must have been a bit wrong! Anyway, pm me with the course that you do. We might already know each other!! 5. I do the Maths Bsc, you? 6. Oops in my original post I forgot to mention that $x^n-y^n=(x-y)\left(\sum^{n-1}_{k=0}{x^ky^{n-1-k}}\right)<0$ only holds if $n \in \mathbb {Z}$ and $n\geq 2$ . Doesn't make much of a difference since the statement is clearly true for $n=1$ . Originally Posted by nmatthies1 The second proof looks really interesting too, but how did you get $(x-y)\left(\sum^{n-1}_{k=0}{x^ky^{n-1-k}}\right)<0$ from $x^n-y^n<0$ ? Without going into a thorough proof...... Say we wanted $x^2-y^2$ . Then we would want $(x-y)(x+y)$ . This is because we want the terms with $x$ and $y$ as a product to cancel. Similarly for $x^3-y^3$ we would want $(x-y)(x^2+xy+y^2)$ and for $x^4-y^4$ we would want $(x-y)(x^3+x^2y+xy^2+y^3)$ . Notice that we always want $(x-y)$ as the first bracket then all terms positive in the next bracket since then the negative terms will cancel out the positive terms because the terms are symmetrical. Hence as a generalisation we can say; $x^n-y^n=(x-y)(x^{n-1}+x^{n-2}y+.....+xy^{n-2}+y^{n-1})$ and thus more formally, $x^n-y^n=(x-y)\left(\sum^{n-1}_{k=0}{x^ky^{n-1-k}}\right)$ Originally Posted by Showcase_22 $\frac{x^k}{y^k}<\frac{y}{x} \Rightarrow \frac{x}{y}<1$ Now the above line doesn't make sense so I thought that the converse couldn't be proven (it states "try" in the question). I must have been a bit wrong! I can't see why that line doesn't make sense...... If $\frac{x^k}{y^k}<\frac{y}{x}$ and $x,y>0$ then it must imply that $x . I think here the direct proof is better. As an added note I might be joining you next year as I have applied to Warwick for a maths degree of course. 7. I like the sum of the series you posted. It's really handy! I'll remember it for a later date. I can't see why that line doesn't make sense...... If and then it must imply that . I think here the direct proof is better. What I was thinking was that since $\frac{y}{x}>1$ then it does not necessarily imply $ \frac{x^k}{y^k}<\frac{y}{x} \Rightarrow \frac{x}{y}<1 $ . My reasoning for this was that $\frac{y}{x}>1$. For example, let's say $\frac{y}{x}=1.2$. Therefore $\frac{x^k}{y^k}<1.2$ so this does not imply $\frac{x}{y}<1$. I would try and prove this using algebra, except when I started to do so it became increasingly confusing to convey my thought processes. I hope using the example of 1.2 made more sense. Anyway, I could be completely wrong so let me know what you think. As an added note I might be joining you next year as I have applied to Warwick for a maths degree of course. You're making us look bad!!! Stop it!!! Seriously, you seem pretty good at maths. What A -levels are you taking and have you applied at Oxford or Cambridge?? Don't think that Warwick is filled with silly people like me. I reckon i'm about average. Some people are (miles) better, some people aren''t so good. I think you'll enjoy it here. 8. Originally Posted by Showcase_22 My reasoning for this was that $\frac{y}{x}>1$. For example, let's say $\frac{y}{x}=1.2$. Therefore $\frac{x^k}{y^k}<1.2$ so this does not imply $\frac{x}{y}<1$. Look at it this way........ say $\frac{y}{x}=1.2 \implies \frac{y^k}{x^k}=1.2^k \implies \frac{x^k}{y^k}=\frac{1}{1.2^k}<1$ and therefore $\frac{x^k}{y^k}<1 \implies \frac{x}{y}<1$ . Does this help? Originally Posted by Showcase_22 Seriously, you seem pretty good at maths. What A -levels are you taking and have you applied at Oxford or Cambridge?? Thank you.... I am taking further maths, economics and physics and yes I have applied to Oxford. However Warwick is a university that is up there right at the top and isn't off the Oxbridge par by much. 9. hmm, I see that..... That's one question I got wrong lol =S Thank you.... I am taking further maths, economics and physics and yes I have applied to Oxford. However Warwick is a university that is up there right at the top and isn't off the Oxbridge par by much. Are you taking STEP and AEA? Some people here have taken both and have S's and merits in the two. If you want an indication of what the Warwick maths homework is like, look for some of my other posts. Most of the stuff I post on here are homework questions I get stuck on.
2014-08-30T15:53:19
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https://brilliant.org/discussions/thread/square-root-algorithm/
× # Square Root Algorithm Fix a number $$\alpha \in \mathbb{R}^+$$. Choose $$x_1 > \sqrt{\alpha}$$ and define: $$x_{n+1} = \frac{1}{2}(x_n +\frac{\alpha}{x_n})$$ We will prove that the sequence of all $$x_n$$: (1) Is bounded below (2)$$\displaystyle \lim_{n \to \infty} x_n = \sqrt{\alpha}$$ (3) Is monotone decreasing Proof: (1) We will show that $$x_n > \sqrt{\alpha}$$ $$\forall$$ $$n \in \mathbb{N}$$. Using induction, we establish a base case as $$n=1$$, which is trivial as this is given: $$x_1 > \sqrt{\alpha}$$ Now, suppose the statement is true for $$n= k-1$$: $$x_{k-1} > \sqrt{\alpha}$$ $$\Rightarrow$$ $$x_{k-1} - \sqrt{\alpha} > 0$$ $$\Rightarrow$$ $$(x_{k-1} - \sqrt{\alpha})^2 > 0$$ $$\Rightarrow$$ $$x_{k-1}^2 -2\sqrt{\alpha}x_{k-1} + \alpha > 0$$ $$\Rightarrow$$ $$x_{k-1}^2 + \alpha > 2\sqrt{\alpha}x_{k-1}$$ $$\Rightarrow$$ $$x_{k-1} +\frac{\alpha}{x_{k-1}} > 2\sqrt{\alpha}$$ $$\Rightarrow$$ $$\frac{1}{2} (x_{k-1} + \frac{\alpha}{x_{k-1}}) > \sqrt{\alpha}$$ $$\Rightarrow$$ $$x_k > \sqrt{\alpha}$$ Then we may conclude that the statement is true for $$n=1$$ and if it is true for some $$n=k-1$$ then it is true for $$n=k$$. The proof follows by induction. This proves (1). (2) Proceeding to take the limit: $$\displaystyle \lim_{n \to \infty} x_{n+1} = \displaystyle \lim_{n \to \infty} \frac{1}{2}(x_n + \frac{\alpha}{x_n})$$ Now, if this limit exists then it must be true that: $$\displaystyle \lim_{n \to \infty} x_n = \displaystyle \lim_{n \to \infty} x_{n+1}$$ Let this limit be represented as $$L$$. Then we have: $$L = \frac{1}{2}(L+ \frac{\alpha}{L})$$ $$\Rightarrow$$ $$2L^2 = L^2 + \alpha$$ $$\Rightarrow$$ $$L = \sqrt{\alpha}$$ Then the sequence is bounded below with $$\displaystyle \lim_{n \to \infty} x_n =\sqrt{\alpha}$$ This proves (2). (3) Using induction, we establish a base case as $$n=1$$: $$\Rightarrow$$ $$x_1 > \frac{1}{2} (x_1 + \frac{\alpha}{x_1})$$ $$x_1 > \sqrt{\alpha}$$ $$\Rightarrow$$ $$\frac{1}{2}x_1^2 > \frac{1}{2} \alpha$$ $$\Rightarrow$$ $$x_1^2 - \frac{1}{2}x_1^2 > \frac{1}{2} \alpha$$ $$\Rightarrow$$ $$x_1^2 > \frac{1}{2} (x_1^2 +\alpha)$$ $$\Rightarrow$$ $$x_1>\frac{1}{2}(x_1 +\frac{\alpha}{x_1})$$ $$\Rightarrow$$ $$x_1>x_2$$ The base case clearly holds. Now the inductive case $$n=k$$: In (1) we proved that $$x_n > \sqrt{\alpha}$$ $$\forall$$ $$n \in \mathbb{N}$$. Begin with: $$x_k > \sqrt{\alpha}$$ $$\Rightarrow$$ $$\frac{1}{2} x_k^2 > \frac{\alpha}{2}$$ $$\Rightarrow$$ $$x_k^2 > \frac{1}{2}x_k^2 + \frac{\alpha}{2}$$ $$\Rightarrow$$ $$x_k > \frac{1}{2}(x_k + \frac{\alpha}{x_k})$$ $$\Rightarrow$$ $$x_k > x_{k+1}$$ The proof again follows by induction. This proves (3). We have shown that the sequence is monotone decreasing and bounded below with a limit at infinity that exists. It follows that the sequence converges. QED Note the utility of this sequence. If we let alpha equal any real number, we now have an algorithm that will approximate the square root of alpha, with accuracy increasing as $$n$$ increases. The proof above is my own, but the sequence itself is not original. I thought this was interesting. Note by Ethan Robinett 2 years, 11 months ago Sort by: To prove (2), you made the assumption that the limit exists. It is not necessarily valid to assume that the limit exists, and then show that the limit is a certain value. You can actually show that the limit exists using (1) and (3). Note that (3) doesn't use (2). That is often the approach taken for this question. Staff · 2 years, 11 months ago I was thinking about this. I read that in order to show that a sequence converges, if it is monotone decreasing, then you must show that it is bounded below, which was done above without (2). So I'm not sure if (2) really needs to be in the proof at all. However, we could have just as easily shown that the sequence is bounded by zero, but the sequence clearly doesn't converge to zero. If (1) and (3) can be used to show that the limit exists, then would (2) be more appropriately placed at the end? · 2 years, 11 months ago The standard approach is to • Show that the sequence is monotonically decreasing. • Show that the sequence is bounded below. • Hence conclude that the limit exists. • Then conclude that the limit is $$\sqrt{\alpha}$$. The order of the first 2 steps is interchangeable. Note that bounded below follows directly from AM-GM. Staff · 2 years, 11 months ago ×
2017-07-24T00:54:46
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https://www.physicsforums.com/threads/why-does-eulers-identity-work-only-in-radians.543866/
Why does Euler's identity work only in Radians? 1. Oct 25, 2011 e^iA = cosA + i*sinA is true iff A is expressed in Radians. Why that particular unit? (I'm not sure this rubric is the right one for this question, but since it didn't seem to fit any of the other rubrics, I put it here.) 2. Oct 25, 2011 D H Staff Emeritus One way to prove Euler's formula is in terms of power series. First express exp(ix) as a power series. Then split the even and odd terms of this to form two series. Next rewrite the powers of i in terms of powers of -1 (and a residual i for the odd series). Those two power series are exactly those of cosine(x) and i*sine(x) where x is in radians. 3. Oct 25, 2011 lurflurf It would work fine with another scaling, except it would have another scaling. given small x we have approximately that exp(i x)~1+i x cos(x)+i sin(x)~1+i x which combined with the addition laws gives us the identity if we had for some reason wanted to use other scaling we have otherexp(i x/a)=exp(i x)~1+i x othercos(x/b)+i othersin(x/b)=cos(x)+i sin(x)~1+i x we would have the identity otherexp(i x/a)=othercos(x/b)+i othersin(x/b) for suitable a and b in calculus we would say a=otherexp'(0) b=othersin'(0) and of course when a=b=1 we have the usual functions Last edited: Oct 25, 2011 4. Oct 25, 2011 Ah, that makes sense. Thank you both, D H and lurflurf 5. Oct 29, 2011 Ah, one detail. In doing the MacLauren series for sin and cos, in the expansion one only needs that sin(0) = 0 and cos(0) = 1. This is true whether or not your units are in degrees or in Radians. Therefore the answer from D H seems to be begging the question. There's something else to this.... What? 6. Oct 29, 2011 Dickfore It all has to do with the following limit: $$\lim_{x \rightarrow 0}{\frac{\sin{x}}{x}} = 1$$ 7. Oct 29, 2011 D H Staff Emeritus Not true. Here is the series for sine(x) when x is in radians: $$\sin x = \sum_{n=0}^{\infty}\frac{(-1)^n}{(2n+1)!}x^{2n+1}$$ This is not valid when x is in degrees. The reason is simple: $d/dx(\sin x) = \cos x$ is only valid when x is in radians. When x is expressed in degrees, $$\frac{d}{dx}\sin x^{\circ} = \frac{\pi}{180} \cos x^{\circ}$$ and thus the Maclaurin series for $\sin x^{\circ}$ is $$\sin x^{\circ} = \sum_{n=0}^{\infty}\frac{(-1)^n}{(2n+1)!}\left(\frac{\pi}{180}\right)^{2n+1}x^{2n+1}$$ 8. Oct 29, 2011 symbolipoint All of those terms are circular functions. Radians are based on pi. 360 Degrees is arbitrary. 9. Oct 29, 2011 HallsofIvy Staff Emeritus You may be thinking that, because (sin x)'= cos x, (sin x)''= -sin x, (sin x)'''= -cos x, etc. and similarly for the derivatives of cos x, in order to find the MacLauren series, we only need to evaluate those at x= 0. But the point is that those derivative formula are only true if x is in radians. In particular, the derivative of sin x is given by $$\lim_{h\to 0}\frac{sin(x+h)- sin(x)}{h}$$ Using the trig identity $sin(x+h)= sin(x)cos(h)+ cos(x)sin(h)$ that becomes $$\lim_{h\to 0}\frac{sin(x)cos(h)- cos(x)sin(h)- sin(x)}{h}$$ $$= \lim_{h\to 0} sin(x)\frac{cos(h)- 1}{h}+ cos(x)\frac{sin(h)}{h}$$ so to have the derivative of sin(x) be cos(x) we must have $$\lim_{h\to 0}\frac{cos(h)- 1}{h}= 0$$ and $$\lim_{h\to 0}\frac{sin(h)}{h}= 1$$ and, as Dickfore noted, those are only true if h is in radian. In fact, if you are going to use the trig functions as general functions and not just as ways to "solve triangles", you should stop thinking of the arguments as being angles at all- and just think of them as number, without any "degrees" or "radians". One way of doing that is to use the "circle" definition: Given a unit circle (circle with center at the origin, radius 1, in some coordinate system), to find sin(t) or cos(t), start from the point with coordinates (1, 0) and measure counter-clockwise (clockwise if t is negative) a distance t around the circumference of the circle. What ever (x,y) point you end at gives you cos(t)= x, sin(t)= y. No angles at all there! (Yes, you have measured a distance but the "units" are whatever "units" you used to construct the coordinate system.) 10. Oct 31, 2011 thanks, HallsofIvy. There are two curious points that I do not quite understand. (1) First, nowhere in your proof is it explicitly shown that x and h need to be in radians. I presume it is implicit in the use of the limit of sin(h)/h = 1 which is derived from using the arc formula arc length = rx, where x is in radians, and r is the radius. Right? (2) Assuming this to be the case, although the proof that sin'(x) = cos x and so forth rests on x being in Radians, nonetheless sin'(x) = cos x works perfectly well for x in degrees, since one just substitutes from the radians to the equivalent number of degrees and back again. The derivative of sin x at x=60 degrees is indeed cos(60 degrees). Therefore I would appreciate some clarification to your statement that these derivatives are only valid if x is in Radians. 11. Oct 31, 2011 lurflurf sin'(0)=1 is only true in radians, it can be taken as the definition of radians. sin'(0)=pi/Arccos(-1) As I said above the sine and cosine in another scaling are easily expressible in terms of the usual one, but that does not make them equal. Some confusion might be resulting from thinking of degrees as a way of expressing a number versus as being a different number. Like if we said 60 degrees=π/3 or degrees=π/180 by sine in degrees we mean a function where (60,√3/2) not (π/3,√3/2)
2017-11-19T05:59:32
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https://math.stackexchange.com/questions/2341887/on-my-attempt-in-showing-that-mathbb-z-sqrt-2-1-epsilon-simeq-math
# On my attempt in showing that $\mathbb Z[\sqrt{-2}] / (1 - \epsilon) \simeq\mathbb Z_3$. Let $A = Z[\sqrt{-2}]$ be the ring with elements of the form $\{ a +b \epsilon\mid a,b \in\mathbb Z \}$ and where $\epsilon^2 = -2$. Let $(1- \epsilon)$ denote the ideal generated by $1- \epsilon$. I am tasked to show that $A / (1 - \epsilon) \simeq \mathbb Z_3$. My attempt: First I see that all elements in the ideal $(1- \epsilon)$ are given by $\{ (a+b\epsilon) (1- \epsilon) | a,b \in Z \}$ so a generic element has the form $a+2b + \epsilon(b -a )$. Now I want to show that there are 3 equivalence classes in $A / (1 - \epsilon)$ where two elements $x,y$ are equivalent $\iff (y_1 +y_2\epsilon) -(x_1+x_2 \epsilon) \in (1 - \epsilon)$. So I obtain the condition that $y_1 - x_1$ must be of the form $a+2b$ and $y_2 -x_2$ must be of the form $b-a$ but I can't seem to find a way to utilize this information. • have you tried solving for $a$ and $b$ in terms of $x_1,x_2,y_1,y_2$ ? – mercio Jun 30 '17 at 13:29 • @mercio yes I tried but I could not get a solution in that way, is it possible? – Monolite Jun 30 '17 at 14:18 Let $u=1- \epsilon$. Then $a+b\epsilon=a+b-bu$. Moreover, $u\bar u=3$. Therefore, $a+b\epsilon \equiv c \bmod u$, where $c=a+b \bmod 3$. This proves that the quotient ring has at most $3$ elements. Since $0,1,-1$ are not equivalent mod $u$, the quotient ring has exactly $3$ elements and so is isomorphic to $\mathbb Z_3$. Alternatively, you can just argue as follows: $$\mathbb Z[\sqrt{-2}] \cong \frac{\mathbb Z[x]}{(x^2+2)} \implies \frac{\mathbb Z[\sqrt{-2}]}{(1-\epsilon)} \cong \frac{\mathbb Z[x]}{(x^2+2,1-x)} = \frac{\mathbb Z[x]}{(3,1-x)} \cong \frac{\mathbb Z}{(3)}$$ using that $x^2+2=(x^2-1)+3=-(1-x)(1+x)+3$. • Thank you greatly, why does the last isomorphism hold? – Monolite Jun 30 '17 at 14:22 • @Monolite, it's induced by $f(x) \mapsto f(1)$. – lhf Jun 30 '17 at 14:34 • There is a theorem (sadly, I don`t remember its proof) that $\mathbb{Z}[i]/(a+bi) \cong \mathbb{Z}_{a^2 + b^2} = \mathbb{Z}_{|a+bi|}$. – Alexander Markov Jun 30 '17 at 14:58 • I thought I understood initially but why do we need that $c=a+b \bmod 3$ what is the link with $u\bar u=3$? – Monolite Jun 30 '17 at 15:33 • @Monolite I posted a complete and rigorous form of the first argument. – Bill Dubuque Jun 30 '17 at 16:29 There is an image of $\,\Bbb Z\,$ in $R = \Bbb Z[\sqrt{-2}]/(1\!-\!\sqrt{-2})\,$ by mapping $\,n\,$ to $\ n \pmod{1\!-\!\sqrt{-2}}$ by composing the natural maps $\,\Bbb Z\to \Bbb Z[\sqrt{-2}]\to R.\,$ This map $\, h\color{#0a0}{ \ {\rm is\ surjective\ (onto)}}$ since $\,{\rm mod}\ \sqrt{-2}\!-\!1\!:\,\ \sqrt{-2}\equiv 1\,\Rightarrow\, a+b\sqrt{-2}\equiv a+b\in\Bbb Z,\,$ and has kernel $\,\color{#c00}{\ker h = 3\Bbb Z}\,$ by \begin{align} n\in (1\!-\!\sqrt{-2})\,\Bbb Z[\sqrt{-2}] \iff &\ \dfrac{n}{1\!-\!\sqrt{-2}}\ \ \,\in\,\ \ \Bbb Z[\sqrt{-2}]\\[0.5em] \iff &\ \dfrac{n(1\!+\!\sqrt{-2})}{3}\in\Bbb Z[\sqrt{-2}] \ \ \ {\rm by\ \ rationalizing\ denominator}\\[0.5em] \iff &\ \color{#c00}{3\mid n}\end{align} By the First Isomorphism Theorem: $\, \color{#0a0}{R = {\rm Im}\ h} \,\cong\, \Bbb Z/\color{#c00}{\ker h} \,=\, \Bbb Z/\color{#c00}{3\Bbb Z}$ • Thank you, was there a more intuitive reason I was missing before? – Monolite Jul 1 '17 at 0:47 From $y_1-x_1 = a+2b$ and $y_2-x_2 = b-a$, you get by adding the two equations, $y_1+y_2-x_1-x_2 = 3b$. If $y_1+y_2-x_1-x_2$ is a multiple of $3$ then your system has a unique solution in integers, $b = (y_1+y_2-x_1-x_2)/3$ and then $a=y_1-x_1-2b$, and so $x_1+\epsilon x_2$ and $y_1+\epsilon y_2$ are in the same class. If $y_1+y_2-x_1-x_2$ is not a multiple of $3$ then the system has a unique solution in the rational numbers but they are not integers so they are not in the same class. Therefore they are in the same class if and only if $x_1+x_2 \equiv y_1+y_2 \pmod 3$. Since $x_1+x_2$ can be congruent to either $0,1$ or $2$ mod $3$, the three equivalence classes are $S_k = \{(x_1+\epsilon x_2) \mid x_1+x_2 \equiv k \pmod 3 \}$ with $k=0,1,2$
2019-09-23T07:01:16
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http://mathhelpforum.com/calculus/147123-word-problem-calculus-print.html
# Word Problem-Calculus • May 31st 2010, 10:40 AM skweres1 Word Problem-Calculus It takes one corn mill 6 minutes to grind a 50 pound bag of corn into cornmeal, while it takes a slower mill 9 minutes to grind the same bag of corn. If both mills are working at the same time, how long would it take to grind 1500 pounds of corn? • May 31st 2010, 10:50 AM Unknown008 If they work together, they should take less time to make a 50 pound bag. This time is given by: $\frac{1}{t_1} + \frac{1}{t_2} = \frac{1}{t_{total}}$ Using the given values, we get: $\frac{1}{6} + \frac{1}{9} = \frac{5}{18}$ So, the total time is 18/5 = 3.6 hours. It takes 3.6 hours to make a 50 pound bag. You can now use proportions to find the time it takes for 1500 pound of corn to be gound. • May 31st 2010, 12:10 PM skweres1 Thanks Thank you, but if it takes 6 minutes for one person to do 50 pounds then how would it take 3.6 hours to to it together? • May 31st 2010, 02:18 PM GeoC Rate Mill 1: $R_1 = 50/6$ Rate Mill 2: $R_2 = 50/9$ Combined Rate $R_t = R_1 + R_2$ Thus: $Time = 1500 / (50/6 + 50/9) = 30/(1/6 + 1/9) = 108 min$ • Jun 1st 2010, 06:25 AM Unknown008 Quote: Originally Posted by skweres1 Thank you, but if it takes 6 minutes for one person to do 50 pounds then how would it take 3.6 hours to to it together? Right, actually, I it was 3.6 minutes, and not 3.6 hours. (Itwasntme) Then, taking from there, you can use proportions to solve the rest. 50 pounds are ground in 3.6 minutes. 1 pound is ground in 3.6/50 minutes. 1500 pounds are ground in 3.6/50 * 1500 = 108 minutes. • Jun 1st 2010, 12:15 PM HallsofIvy Quote: Originally Posted by skweres1 It takes one corn mill 6 minutes to grind a 50 pound bag of corn into cornmeal, while it takes a slower mill 9 minutes to grind the same bag of corn. If both mills are working at the same time, how long would it take to grind 1500 pounds of corn?
2017-08-18T12:56:12
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https://www.physicsforums.com/threads/square-roots-of-large-numbers.555836/
# Square roots of large numbers 1. Dec 1, 2011 ### Quark Itself I stumbled upon a math question, at the glance of it, seemed easy. One is supposed to find the exact value of this square root: √30*31*32*33+1 They are all under the square root operator and Fundamental BEDMAS/BIDMAS applies of course. The trick here is when the condition states that one cannot use a calculator. Any suggestions, anyone? 2. Dec 1, 2011 ### jgens I do not know any tricks to actually solve this, but you can get a fairly good estimate using basic calculus. In particular, note that 312 < (30*31*32*33 + 1)1/2 < 322. We can model √ around 314 linearly by the function defined by L(x) = 312 + (2*312)-1(x-314). Take x = 31.54 as an initial guess and we get something like L(x) ≈ 992. I think that actual answer is 991, so this gets you a pretty good estimate and all of the calculations are manageable. I am sure someone else knows a nice method to tackle this problem though, so I would like to see that as well. 3. Dec 1, 2011 ### D H Staff Emeritus You are looking for some number $n=\sqrt{30*31*32*33+1}$. Squaring both sides to get rid of the radical yields $n^2=30*31*32*33+1$. Writing 31, 32, and 33 as 30+1, 30+2, and 30+3 yields $n^2=30^4+6*30^3+11*30^2+6*30+1$. This suggests writing n in a similar form: $n=30^2+a*30+b$. Squaring yields $n^2 = 30^4+2a*30^3+(a^2+2b)*30^2+2ab*30+b^2$. Equating the expansion of 30*31*32*33+1 with this readily yields a=3, b=1, and thus n=991. 4. Dec 1, 2011 ### D H Staff Emeritus This is another approach. The answer has to be fairly close to 992. The word "exact" in the question indicates the number has to be an integer. The last digit of this number must be either one or nine, and the answer certainly isn't 999 or 981. The only viable choices are 991 and 989. Just square each (not hard) and compare to 30*31*32*33+1 (which you will have to expand for this approach). Since the original guess was 992, starting with 991 will eliminate the need to try 989 as the second choice since the answer is indeed 991. 5. Dec 1, 2011 ### jgens That is a nice solution! Much better than my thoughts on the problem. 6. Dec 1, 2011 ### Quark Itself Thanks for the quick replies ! Really appreciate it ! I was thinking along the lines of having a number n which would be that exact value but as I worked it through I got stuck, where you presented the n = 30^2 + 30a + b part was where I got stuck, but your explanation was really good so thanks again. 7. Dec 2, 2011 ### lurflurf 30*31*32*33+1=x^2 since 31*32-30*33=(x+1)-(x-1) and 30*31*32*33=(x+1)*(x-1) we can see that 30*33=x-1 31*32=x+1 and x is easy to find
2018-05-24T00:40:42
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https://www.physicsforums.com/threads/what-does-dimensionally-mean.41582/
# Homework Help: What does dimensionally mean? 1. Sep 3, 2004 ### tony873004 Here's a homework question I have. I don't even understand what they're asking. Does somebody know the answer, and can you explain it in a way that I can do the other problems like it?? Thanks in advance for your help :rofl: Velocity is related to acceleration and time by the following expression, v = a^p * t^q. Find the powers p and q that make this equation dimensionally consistent. 2. Sep 3, 2004 ### enigma Staff Emeritus Hi tony, In the future, questions like this should be posted in the homework help forums. What the question is asking you is to try to make the units multiply out correctly. It is one of the easiest ways to ensure that you've done a problem correctly. Instead of using numbers to multiply together, you want to multiply the units together. You know that velocity has units of length/time. You know that acceleration has units of length/time^2. If you treat the units as variables, you should notice what values of p and q which make the equality sign true. Give it a shot... write back if you're still stuck. 3. Sep 3, 2004 ### tony873004 I just figured that out! Sorry, it was my first post The left side of the = is L / T^2. Does this mean that I have to make the right side equal L / T^2 ? If that's the case, then (L / T) should have a power of one, and T should have whatever power allows me to combine it with the T in the L/T denominator to get T^2. L / T^2 = (L / T) ^ p * T ^q I asked the teacher this question after class today, he looked at it quickly and said that p =1 and q = -1. But we do our homework on Webassign.net, which instantly tells you if your answers are correct. The p=1 was correct, but the q=-1 wasn't. Last time I ask my teacher for help I guess I just forget the algebra that lets me move the T from outside the parenthesis to the denominator of the expression inside the parenthesis. (p.s. I saw on someone else's post that there's a nicer way to say stuff like T^q. A way that makes it look more textbook by actually placing the q as a superscript of T. How did he do that?) 4. Sep 3, 2004 ### Grizzlycomet $$T^q$$ Its called Latex. Check out this page to see how to make it :rofl: Oh, and welcome to PF! 5. Sep 3, 2004 ### Mk Yes! Welcome! 6. Sep 3, 2004 ### JasonRox Sometimes even if you forget the equations, you can get the answer simply by looking at the dimensions, and working it out. Of course it has to make sense, but it is a good step towards understanding. 7. Sep 3, 2004 ### Galileo $$v=a^pt^q$$ Velocity $v$ is measured in meters per second (m/s). Acceleration $a$ is measured in meters per second squared $(m/s^2)$. Time $t$ is ofcourse measured in seconds (s). So the left side is m/s. 8. Sep 3, 2004 ### arildno Let's look at how you can directly compute p and q using "unit maths" $$[v]=L^{1}T^{-1}$$ $$[a]=L^{1}T^{-2}$$ $$[t]=T^{1}$$ Hence, we require: $$[v]=[a^{p}t^{q}]=[a^{p}][t^{q}]=[a]^{p}[t]^{q}$$ Or, filling in: $$L^{1}T^{-1}=L^{1*p}T^{-2*p}T^{1*q}=L^{p}T^{q-2p}$$ Now, L and T are independent dimensions, so the powers must equal each other seperately: L: 1=p T: -1=q-2p Hence, you may solve for: p=1, and q=1 9. Sep 3, 2004 ### HallsofIvy It does not follow that the teacher "screwed up". In his post, tony873004 asserted that "The left side of the = is L / T^2" which was incorrect. If he gave that to his teacher when he was asking the problem, then his teacher was correct. I suspect he didn't show his teacher the original problem. 10. Sep 3, 2004 ### HallsofIvy tony873004: when you are asking us or your teacher for help (and you should certainly continue doing both) be sure to state the ENTIRE PROBLEM and show what you have already done. 11. Sep 3, 2004 ### tony873004 Arildno's explanation made sense to me. It's still going to take a little time to digest this, but p=1 and q=1 are correct according to Webassign.net. I did post the entire question. The $$L / T^2$$ was my screw-up, not the teacher's, but I made this mistake after talking to the teacher. He did see the assignment straight from the book. Thanks everyone for helping. This is so awesome that there's a forum of people willing to help!! Now I've got a $$26/26$$ on this assignment . And TEX is very cool
2018-08-17T04:03:55
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https://www.physicsforums.com/threads/row-reduction-to-solve-for-6-unknowns.747767/
# Row Reduction to solve for 6 Unknowns 1. Apr 8, 2014 ### Yosty22 1. The problem statement, all variables and given/known data I am in a calculus class where we are learning the introduction to row reduction. I have done this before in other courses, so I am familiar with the process, but I am not sure about this one. We were given: x4 + 2x5- x6 = 2 x1 + 2x2 + x5 -x6 = 0 x1 + 2x2 + 2x3 - x5 + x6 = 2 We were supposed to: solve for each unknown or tell how many solutions this has(usually something like 0 or infinity). 2. Relevant equations 3. The attempt at a solution We haven't learned it exactly yet, but there is no real row eschelon form for this, is there? I don't know what all I can do mathematically, it just seems to me that we only have 3 equations and 6 unknowns, so we cannot possibly solve for this, right? Am I missing something here? I wrote out the matrix exactly as it is written (if it didn't state a variable, i set it to 0. e.g. in equation 1, it does not mention x1 through x3 so I made those 0). No matter what you do, you have a 3x6 matrix with 3 equations and 6 unknowns. Am I missing something here? How could I mathematically answer this question and not just state what I have stated above? Any help would be appreciated, thanks. 2. Apr 8, 2014 ### LCKurtz Just go ahead and row reduce it as far as you can. You can get rid of the $x_1$ in the middle equation with the last one, which will also get rid of the $x_2$ but that's OK. You should be able to get three rows whose first element is $1$. Solve for those variables in terms of the others. 3. Apr 8, 2014 ### Yosty22 That's what I thought, thanks. So I have reduced it as far as I can, and it is obvious that you still cannot solve, so would an appropriate answer just be that there are No Solutions? 4. Apr 8, 2014 ### LCKurtz I'm thinking you can likely solve for three of the variables in terms of the others, right? The extra variables are called free variables and can be anything. You have lots of solutions. The case when there would be no solutions is of one of your rows is all zeros and a nonzero on the right side. 5. Apr 9, 2014 ### HallsofIvy Staff Emeritus What do you mean by "cannot solve"? You cannot solve for a single specific value for each number but you didn't expect to, right? Part of the question was "how many solutions are there?" You have three equations in six unknown values so if all three equations are independent you would expect to be able to solve for 6- 3= 3 of the values in terms of the other three.
2017-11-20T12:38:19
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https://uvacs2102.github.io/class16/
Discrete Mathematics Class 16: Structural Induction Schedule Problem Set 6 is due tomorrow at 6:29pm. Make sure to read the corrected version of Problem 7. Lists Definition. A list is an ordered sequence of objects. A list is either the empty list ($\lambda$), or the result of $\text{prepend}(e, l)$ for some object $e$ and list $l$. \begin{equation} \begin{split} \text{\em first}(\text{prepend}(e, l)) &= e \text{\em rest}(\text{prepend}(e, l)) &= l \text{\em empty}(\text{prepend}(e, l)) &= \text{\bf False} \text{\em empty}(\text{\bf null}) &= \text{\bf True} \end{split} \end{equation } Definition. The length of a list, $p$, is: \begin{equation} \begin{split} \begin{cases} 0 & \text{if}\ p\ \text{is \bf null} \text{\em length}(q) + 1 & \text{otherwise}\ p = \text{\em prepend}(e, q)\ \text{for some object}\ e\ \text{and some list}\ q \end{cases} \end{split} \end{equation } def list_length(l): if list_empty(l): return 0 else: return 1 + list_length(list_rest(l)) Prove: for all lists, $p$, list_length(p) returns the length of the list $p$. Concatenation Definition. The concatenation of two lists, $p = (p_1, p_2, \cdots, p_n)$ and $q = (q_1, q_2, \cdots, q_m)$ is $$(p_1, p_2, \cdots, p_n, q_1, q_2, \cdots, q_m).$$ Provide a constructuve definition of concatenation. Note that $\text{prepend}(p,q)$ is not a good idea for two reasons. If we use this definition, then the first element of the constructed list will be the object (list) $p$ (as a whole) rather than the first element $p_1$ of the list $p$. Also, if we want to only define lists of specific objects, for example integers, we can still use the same recursive/constructive definition of lists by substituting “object” with “integer”, but in that case $\text{prepend}(p,q)$ will not even well defined, as it can only accept integers as first input. # Structural Induction To prove proposition $P(x)$ for element $x \in D$ where $D$ is a recursively-constructed data type, we do two things: 1. Show $P(x)$ is true for all $x \in D$ that are defined using base cases. 2. Show that if $P(y)$ is true for element $y$ and $x$ is constructed from $y$ using any “construct case” rules, then $P(x)$ is true as well. Comparing Various forms of Induction \begin{center} \begin{tabular}{lccc} & {\bf Regular Induction} & {\bf Invariant Principle} & {\bf Structural Induction} \ \hline Works on: & natural numbers & state machines & data types To prove $P(\cdot)$ & {\em for all natural numbers} & {\em for all reachable states} & {\em for all data type objects} Prove {\bf base case(s)} & $P(0)$ & $P(q_0)$ & $P(\text{base object(s)})$ and {\bf inductive step} & $\forall m \in \mathbb{N} \ldotp$ & $\forall (q, r) \in G \ldotp$ & $\forall s \in \text{\em Type} \ldotp$ & $P(m) \implies P(m+1)$ & $P(q) \implies P®$ & $P(s) \implies P(t)$ & & & $\quad \forall t\ \text{constructable from}\ s$ \ \end{tabular} \end{center} # ## Prove. For any two lists, $p$ and $q$, $\text{length}(p + q) = \text{length}(p) + \text{length}(q)$. #
2019-12-16T04:30:54
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https://math.stackexchange.com/questions/3159746/what-is-the-difference-between-cauchys-theorem-and-cauchys-theorem-for-abelian
# What is the difference between Cauchy's Theorem and Cauchy's Theorem for Abelian Groups? These theorems seem to be identical but for some reason, the requirement that a group is finite AND abelian is sometimes stated instead of just finite. Could someone let me know if there is a difference in these theorems? ## 2 Answers Cauchy's theorem is for all finite groups not just for abelians. The abelian case is somehow easier to be proved. All in all, if something is true for all groups then it is true for abelian groups whatsoever. Cauchy's Theorem is for all finite groups. The proof is sometimes done separately for the abelian case and the non-abelian case. However this is not necessary, see here: Cauchy's Theorem for Abelian Groups • It's just easier to prove the latter, especially in introduction to algebra texts. – Don Thousand Mar 23 at 20:05
2019-04-21T02:42:41
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http://prairiepointpta.org/gmhdlb7v/what-is-bijective-function-72477a
# what is bijective function As pointed out by M. Winter, the converse is not true. My examples have just a few values, but functions usually work on sets with infinitely many elements. In essence, injective means that unequal elements in A always get sent to unequal elements in B. Surjective means that every element of B has an arrow pointing to it, that is, it equals f(a) for some a in the domain of f. A function f : A -> B is said to be onto function if the range of f is equal to the co-domain of f. How to Prove a Function is Bijective without Using Arrow Diagram ? A function that is both One to One and Onto is called Bijective function. The function f is called as one to one and onto or a bijective function, if f is both a one to one and an onto function. And I can write such that, like that. Mathematical Functions in Python - Special Functions and Constants; Difference between regular functions and arrow functions in JavaScript; Python startswith() and endswidth() functions; Hash Functions and Hash Tables; Python maketrans() and translate() functions; Date and Time Functions in DBMS; Ceil and floor functions in C++ $$Now this function is bijective and can be inverted. If it crosses more than once it is still a valid curve, but is not a function. A function is invertible if and only if it is a bijection. Thus, if you tell me that a function is bijective, I know that every element in B is “hit” by some element in A (due to surjectivity), and that it is “hit” by only one element in A (due to injectivity). Question 1 : The figure shown below represents a one to one and onto or bijective function. Definition: A function is bijective if it is both injective and surjective. Stated in concise mathematical notation, a function f: X → Y is bijective if and only if it satisfies the condition for every y in Y there is a unique x in X with y = f(x). Some types of functions have stricter rules, to find out more you can read Injective, Surjective and Bijective. So we can calculate the range of the sine function, namely the interval [-1, 1], and then define a third function:$$ \sin^*: \big[-\frac{\pi}{2}, \frac{\pi}{2}\big] \to [-1, 1]. Each value of the output set is connected to the input set, and each output value is connected to only one input value. Ah!...The beautiful invertable functions... Today we present... ta ta ta taaaann....the bijective functions! The inverse is conventionally called $\arcsin$. More clearly, f maps distinct elements of A into distinct images in B and every element in B is an image of some element in A. Hence every bijection is invertible. Infinitely Many. This is equivalent to the following statement: for every element b in the codomain B, there is exactly one element a in the domain A such that f(a)=b.Another name for bijection is 1-1 correspondence (read "one-to-one correspondence).. A bijective function is both injective and surjective, thus it is (at the very least) injective. Functions that have inverse functions are said to be invertible. Below is a visual description of Definition 12.4. In mathematics, a bijective function or bijection is a function f : A → B that is both an injection and a surjection. And a function is surjective or onto, if for every element in your co-domain-- so let me write it this way, if for every, let's say y, that is a member of my co-domain, there exists-- that's the little shorthand notation for exists --there exists at least one x that's a member of x, such that. A valid curve, but functions usually work on sets with infinitely many elements both! 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Just a few values, but functions usually work on sets with infinitely many.!
2021-07-30T05:28:56
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https://math.stackexchange.com/questions/78326/expressing-the-venn-diagram/78334
# Expressing the Venn diagram I need help with this question, Express the shaded portion of the Venn diagram in terms of unions, intersections or complements of sets. I have done it as A Intersection B + A Intersection C - A Intersection B Intersection C + D Intersection E + (E-C) i.e. $A\cap B \cup A\cap C - A\cap B\cap C \cup D \cap E \cup (E-C)$ • $(A\cap B\cap C)$ is counted in $(A\cap B) + (A\cap C)$ twice, so if you subtract $(A\cap B\cap C)$, you're still left with one. – user13888 Nov 2 '11 at 21:16 • @MikeWierzbicki: I just missed that. Thanks – user2857 Nov 2 '11 at 21:17 • Also be careful with what you're taking out when you add $(E-C)$. – user13888 Nov 2 '11 at 21:17 • @Akito: I've added TeX-ed version $A\cap B \cup A\cap C - A\cap B\cap C \cup D \cap E \cup (E-C)$ to make post more readable. Maybe you might have a look if this captures your intended meaning. You might also think about adding some parenthesis there, which might improve readability too. (I've tried to stick to your original text, just replaced $+$ by $\cup$ and intersection by $\cap$. – Martin Sleziak Nov 11 '11 at 8:46 You are asked to express the sets in terms of union, intersection, complements. In particular, the "minus" in your proposed solution would have to be expressed in terms of these. But aside from that minor comment, your answer is correct. We do the problem in another way. We look at the four shaded parts, one after the other, going roughly downwards. The top shaded part is inside $A$, inside $B$, outside $C$, outside $D$, and outside $E$. In symbols, it is $$A\cap B\cap C^c\cap D^c\cap E^c,$$ ($S^c$ denotes the complement of $S$). The next one down is very similar, except we are outside $B$, inside $C$, with the rest exactly as above: $$A\cap B^c\cap C\cap D^c\cap E^c.$$ The next one down is outside $A$, outside $B$, and inside $C$, $D$, and $E$: $$A^c\cap B^c\cap C\cap D\cap E.$$ Finally, the one at the bottom is inside $E$, and outside all the others: $$A^c\cap B^c\cap C^c\cap D^c\cap E.$$ Finally, we take the union of our four pieces. It is long, so we use a smaller font: $$\small(A\cap B\cap C^c\cap D^c\cap E^c)\cup(A\cap B^c\cap C\cap D^c\cap E^c)\cup(A^c\cap B^c\cap C\cap D\cap E)\cup(A^c\cap B^c\cap C^c\cap D^c\cap E).$$ Comment: This was a bit lengthy, but utterly mechanical. There are many other forms for the answer. For the purpose of circuit design, we would almost certainly want an expression that used fewer symbols. (We used a total of $19$ $\cap$ and/or $\cup$. The complement operations were not counted, because in circuit design they are often assumed to be "for free.") • $(A\cap B\cap C^c)\cup(A\cap B^c\cap C)\cup( D\cap E)\cup( C^c\cap E)$ would be shorter though not as systematic – Henry Nov 2 '11 at 21:48 • @Henry: Thank you, the comment will be helpful to the OP. I wanted to give a procedure that was as mechanical as possible. – André Nicolas Nov 2 '11 at 22:13 • Thankyou very good answer. – user2857 Nov 3 '11 at 5:53 Note that there are several ways to express the shaded area. Let us go through the parts: • There's a part of $A\cap B$, and a part of $A\cap C$. What missing from both is $A\cap B\cap C$. So we can write this part as $A\cap(B\triangle C)$, where $B\triangle C$ is the symmetric difference, that is $(B\cup C)\setminus(B\cap C)$. • On the bottom part there is a part of $E$ which meets $D$, and the part of $E$ which is not in $C$. We can write this as $(E\cap D)\cup(E\setminus C)$ Together we can write it all as: $$(A\cap(B\triangle C))\cup(E\cap D)\cup(E\setminus C)$$ • I also wanted to ask if an area is not shaded, can we subtract it twice or more times? Does it matter? – user2857 Nov 2 '11 at 21:19 • @Akito: The definition is: $$A\setminus B=\{x\mid x\in A\land x\notin B\}$$ Now suppose what is $(A\setminus B)\setminus B$? These are all the elements that are in $A\setminus B$ and not in $B$. However all the elements of $A\setminus B$ are not in $B$ by definition. So it does not matter to subtract twice. – Asaf Karagila Nov 2 '11 at 21:58
2020-02-16T22:45:29
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https://math.stackexchange.com/questions/2709325/need-help-with-simple-mathematical-sets-difference-intersection-problem
# Need Help With Simple Mathematical Sets Difference (Intersection) Problem I'm having trouble with a seemingly simple problem in Math and I need some help with it. The problem states: A hospital is doing a treatment research on 50 volunteers. Of those, some have had a few reactions: 12 have had headaches, 8 felt nausea and 4 had headaches and nausea. How many volunteers had headaches and not nausea? How many volunteers didn't feel neither (headache and nausea simultaneously)? It's not a homework as I'm not in school anymore but studying by myself to try my country's equivalent of the SAT to get into college. So, I know that my Universe (U) in this case is 50. I know the sets (let's say A and B) are A={12}, B={8} and A∩B={4}, so the number of volunteers that had only headaches exclusively is 8 (12 - 4 people that also had nausea) but i can't figure out the number of volunteers that haven't had neither. I've tried putting together a simple equation to find the number: 12 volunteers that felt headache + 8 that felt nausea + x people that felt neither = 50 that will amount to x=30 volunteers but the book (without explaining why) says the correct answer is 34. I've tried breaking the problem down by its inquiries, writing them one by one (which I won't do here to not extend the question) and it seems to be a simple thing, really, and that what I've done is correct. Is the book wrong (unlikely)? Is my solution correct? How should I go about when solving this problem? • Have you tried drawing a Venn diagram? How many had nausea but not headache? – saulspatz Mar 26 '18 at 19:36 • Yes but it's just like the numbers I've stated on the question, the diagram represents the same thing. It doesn't give me any information that I don't know. – Pedro _PR Mar 26 '18 at 19:40 • Can you find how many distinct volunteers had either headaches or nausea or both from what you've already done? – K B Dave Mar 26 '18 at 19:41 • @Pedro_PR If you fill out the numbers in the Venn diagram carefully, you'll see what's missing. – saulspatz Mar 26 '18 at 19:46 headache and no nausea = 12 - 4 nausea and no heaqdache = 8 - 4 neither = 50 - (12 + 8 - 4) = 34 HINT There are four mutually exclusive categories of volunteers: The sum of the numbers in these groups must add up to $50.$ Are you certain you've counted each group correctly? (No, you haven't, actually, although your calculation of those with headache but not nausea is correct.)
2021-07-26T21:05:12
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https://math.stackexchange.com/questions/2842995/is-there-a-way-to-show-that-ex-lim-n-to-infty-left1-fracxn-rightn
# Is there a way to show that $e^x=\lim_{n\to \infty }\left(1+\frac{x}{n}\right)^n$? I know that $$e:=\lim_{n\to \infty }\left(1+\frac{1}{n}\right)^n,$$ by definition. Knowing that, I proved successively that $$e^{k}=\lim_{n\to \infty }\left(1+\frac{k}{n}\right)^n,$$ when $k\in \mathbb N$, $k\in \mathbb Z$ and $k\in\mathbb Q$. Now, I was wondering : how can I extend this result over $\mathbb R$ ? I tried to prove that $f_n(x):=(1+\frac{x}{n})^n$ converge uniformly on $\mathbb R$ but unfortunately it failed (I'm not sure that it's even true). Any idea ? My idea was to define the function $x\longmapsto e^x$ as $$e^x=\begin{cases}e^x& x\in \mathbb Q\\ \lim_{n\to \infty }e^{k_n}&\text{if }k_n\to x \text{ and }(k_n)\subset \mathbb Q\end{cases}.$$ But to conclude that $$e^x=\lim_{n\to \infty }\left(1+\frac{x}{n}\right)^n,$$ I need to prove that $f_n(x)=\left(1+\frac{x}{n}\right)^n$ converge uniformly on a neighborhood of $x$, but I can't do it. I set $$g_n(x)=f_n(x)-e^x,$$ but I can't find the maximum on a compact that contain $x$, and thus can't conclude. • How do you define $\exp(x)$ for an arbitrary $x$, by the way? Anyhow, there exists a sequence $\left(x_n\right)_{n\in\mathbb{Z}_{>0}}$ of rational numbers such that $x_n\to x$ as $n\to\infty$. – Batominovski Jul 6 '18 at 16:16 • as $\sum_{k=0}^\infty \frac{x^k}{k!}$ – MathBeginner Jul 6 '18 at 16:20 • See this. Not sure it's an exact enough duplicate? – Jyrki Lahtonen Jul 6 '18 at 16:34 • math.stackexchange.com/questions/1557074/… – Robert Wolfe Jul 6 '18 at 16:38 • You may also try to prove this for complex $x=a+ib$ given the definition $e^{a+ib} =e^a(\cos b+i\sin b)$. – Paramanand Singh Jul 6 '18 at 19:18 We can use that exists $p_n, q_n \in \mathbb{Q}$ such that $p_n,q_n \to x$ and $p_n\le x\le q_n$, therefore $$\left(1+\frac{p_n}{n}\right)^n\le \left(1+\frac{x}{n}\right)^n\le \left(1+\frac{q_n}{n}\right)^n$$ and $$\left(1+\frac{p_n}{n}\right)^n=\left[\left(1+\frac{p_n}{n}\right)^\frac{n}{p_n}\right]^{p_n}\to e^x$$ $$\left(1+\frac{q_n}{n}\right)^n=\left[\left(1+\frac{q_n}{n}\right)^\frac{n}{q_n}\right]^{q_n}\to e^x$$ indeed for $\frac{n}{p_n}\in (m,m+1)$ with $m\in \mathbb{N}$ we have $$\left(1+\frac1{m+1}\right)^m\le \left(1+\frac{p_n}{n}\right)^\frac{n}{p_n}\le \left(1+\frac1m\right)^{m+1}$$ and therefore $\left(1+\frac{p_n}{n}\right)^\frac{n}{p_n}\to e$. • I think this is by far the simplest approach, and it does not rely on other properties of the exponential (e.g. its derivative or the properties of the logarithm as its inverse). The key observation is basically that $x$ can be moved from the inside of the parentheses to the exponent by substituting the variable used in the limit. – sasquires Jul 6 '18 at 16:35 • Sorry but why $\left(\left(1+\frac{p_n}{n}\right)^{\frac{p_n}{n}}\right)^{p_n}\to e^x$ ? – MathBeginner Jul 6 '18 at 16:37 • @MathBeginner since $p_n/n \to 0$ we have $$\left(1+\frac{p_n}{n}\right)^{\color{red}{\frac{n}{p_n}}}\to e$$ and $p_n \to x$. – gimusi Jul 6 '18 at 16:39 • Since the exponent depend on $n$, we can't do that, no ? – MathBeginner Jul 6 '18 at 16:43 • @MathBeginner Yes we can since $p_n/n \to 0$. – gimusi Jul 6 '18 at 17:12 It is not difficult to prove the result for real irrational $x$ if you have already proved the case for rational $x$. The only idea you need to establish first as a part of your definition of $e^x$ is that $f(x) =e^x$ is continuous everywhere. I leave this as an exercise for you (hint: show that $\lim_{x\to 0}e^x=1$ using your definition). Now let $x$ be any irrational number. Given any $\epsilon>0$ there is $\delta>0$ such that $$e^x-\epsilon<e^t<e^x+\epsilon$$ whenever $|t-x|<\delta$. Consider two rationals $r, s$ with $x-\delta<r<x<s<x+\delta$ and then we have $$e^x-\epsilon <e^r<e^s<e^x+\epsilon$$ Now we have $$\left(1+\frac{r}{n}\right)^n<\left(1+\frac{x}{n}\right)^n<\left(1+\frac{s}{n}\right)^n$$ and taking limits as $n\to\infty$ we get $$e^x-\epsilon<e^r\leq \lim_{n\to\infty} \left(1+\frac{x}{n}\right)^n\leq e^s<e^x+\epsilon$$ (the above assumes that the limit in question exists for irrational $x$ also and you can prove it using the fact that a bounded monotone sequence is convergent, or better apply liminf/limsup to the above inequalities). Since $\epsilon$ is arbitrary it follows that $$e^x=\lim_{n\to\infty} \left(1+\frac{x}{n}\right)^n$$ Based on feedback from Mark Viola via comments I am giving a link to my blog posts which discuss various routes to the theory of exponential and logarithmic functions : • (+1) Well done, as always! – Mark Viola Jul 6 '18 at 17:24 • @MarkViola: thanks. Btw this exponential and logarithmic functions of a real variable is one of my favorite topics. – Paramanand Singh Jul 6 '18 at 17:26 • Yes, I know. I have read your "blog" and its three-part series. It is well-written. I suggest that you consider embedding a link to it as a solid reference on this topic. – Mark Viola Jul 6 '18 at 17:27 • @mathbeginner This is the answer that you covet. – Mark Viola Jul 6 '18 at 17:31 • @MarkViola: added the links. – Paramanand Singh Jul 6 '18 at 17:33 To prove $$e^x=\lim_{n\to \infty }\left(1+\frac{x}{n}\right)^n$$ Let $$y=\left(1+\frac{x}{n}\right)^n$$ $$\ln y=n \ln(1+x/n)$$ $$= \frac {\ln(1+x/n)}{(1/n)}$$ $$\lim_{n\to \infty }\ln y=\lim_{n\to \infty }\frac {\ln(1+x/n)}{(1/n)}=x$$ Thus $$\lim_{n\to \infty } y= e^x$$ • This doesn't address the question, which is "How does one show that $$\left(\lim_{n\to\infty}\left(1+\frac1n\right)^n\right)^x =\lim_{n\to \infty}\left(1+\frac xn\right)^n?"$$ – Mark Viola Jul 6 '18 at 16:29 $$\frac xn(\frac n{n+x})\le\int_1^{1+\frac xn}\frac1t dt\le\frac xn(1)\implies \frac x{n+x}\le\ln (1+\frac xn)\le\frac xn\implies e^{\frac x{n+x}}\le1+\frac xn\le e^{\frac xn}\implies e^{\frac{xn}{n+x}}\le(1+\frac xn)^n\le e^x\implies e^x\le\lim_{n\to\infty}(1+\frac xn)^n\le e^x$$, by the squeeze or sandwich theorem... Left side: The exponential function may be written as a Taylor series: $e^x=1+x+\frac{x^2}{2!}+\frac{x^3}{3!}+...$ Right side: $(1+\frac{x}{n})^n$ is a binomial expansion like: $(1+y)^n=\binom{n}{0}y^0+\binom{n}{1}y^1+\binom{n}{2}y^2+...+\binom{n}{n-1}y^{n-1}+\binom{n}{n}y^n$ Where $\binom{n}{k}$ is the Binomial coefficient given by the formula : $\binom{n}{k}=\frac{n!}{k!(n-k)!}$ Some basic properties of $\binom{n}{k}$: a)$\binom{n}{0}=1$ because $\frac{n!}{0!(n-0)!}=\frac{n!}{1*n!}$ b)$\binom{n}{1}=n$ because $\frac{n!}{1!(n-1)!}=\frac{(n-1)!*n}{(n-1)!}$ c)$\binom{n}{n-1}=n$ because $\frac{n!}{(n-1)!(n-(n-1))!}=\frac{(n-1)!*n}{(n-1)!*1!}$ d)$\binom{n}{n}=1$ because $\frac{n!}{n!(n-n)!}=\frac{1}{1!}$ e) The formula does exhibit a symmetry that is less evident from the multiplicative formula: $\binom{n}{k}=\binom{n}{n-k}$ Returning: $(1+\frac{x}{n})^n=1+n*\frac{x}{n}+\frac{n!}{2!(n-2)!}\frac{x^2}{n^2}+\frac{n!}{3!(n-3)!}\frac{x^3}{n^3}+...+\frac{n!}{3!(n-3)!}\frac{x^{n-3}}{n^{n-3}}+\frac{n!}{2!(n-2)!}\frac{x^{n-2}}{n^{n-2}}+n*\frac{x^{n-1}}{n^{n-1}}+\frac{x^n}{n^n}$ $(1+\frac{x}{n})^n=1+x+\frac{(n-1)n}{n^2}\frac{x^2}{2!}+\frac{(n-2)(n-1)n}{n^3}\frac{x^3}{3!}+...+\frac{(n-2)(n-1)n}{3!}\frac{x^{n-3}}{n^{n-3}}+\frac{(n-1)n}{2!}\frac{x^{n-2}}{n^{n-2}}+\frac{x^{n-1}}{n^{n-2}}+\frac{x^n}{n^n}$ $(1+\frac{x}{n})^n=1+x+\frac{n-1}{n}\frac{x^2}{2!}+\frac{(n-2)(n-1)}{n^2}\frac{x^3}{3!}+...+\frac{(n-2)(n-1)}{n^{n-4}}\frac{x^{n-3}}{3!}+\frac{n-1}{n^{n-3}}\frac{x^{n-2}}{2!}+\frac{x^{n-1}}{n^{n-2}}+\frac{x^n}{n^n}$ Let's analyze what happens for $n\rightarrow\infty$-here we have three types of limits: -First type: $\displaystyle\lim_{n \to \infty}\frac{n-1}{n}=\displaystyle\lim_{n \to \infty}[1+\frac{1}{n}]=1+0=1$ $\displaystyle\lim_{n \to \infty}\frac{(n-2)(n-1)}{n^2}=\displaystyle\lim_{n \to \infty}\frac{n^2-3n+2}{n^2}=\displaystyle\lim_{n \to \infty}[1-\frac{3}{n}+\frac{2}{n^2}]=1-0+0=1$ Hense $\displaystyle\lim_{n \to \infty}\Bigg(\frac{\displaystyle\prod_{i=1}^{k} (n-i)}{n^k}\Bigg)=1$ -Second type is $\displaystyle\lim_{n \to \infty} \frac{x^{n-\alpha}}{n^{n-\beta}}$-Because ${n^{n-\beta}}$ grows much fasten than $x^{n-\alpha}$ hense: $\displaystyle\lim_{n \to \infty} \frac{x^{n-\alpha}}{n^{n-\beta}}=0$ -Third type: $\displaystyle\lim_{n \to \infty}\Bigg(\frac{\displaystyle\prod_{i=1}^{k} (n-i)}{n^{n-k-1}}\frac{x^{n-k}}{k!}\Bigg)$ We have to show on the biggest power (similar to the first type) as the most relevant: $\frac{\displaystyle\prod_{i=1}^{k} (n-i)}{n^{n-k-1}}\frac{x^{n-k}}{k!}\sim\frac{n^{k-1} }{n^{n-k-1}}\frac{x^{n-k}}{k!}=n^{k-1-(n-k-1)}\frac{x^{n-k}}{k!}=n^{2k-n}*\frac{x^{n-k}}{k!}=\frac{1}{k!}*\frac{x^{n-k}}{n^{n-2k}}$ Again: ${n^{n-\beta}}$ grows much faster than $x^{n-\alpha}$ Hense: $\displaystyle\lim_{n \to \infty}\Bigg(\frac{\displaystyle\prod_{i=1}^{k} (n-i)}{n^{n-k-1}}\frac{x^{n-k}}{k!}\Bigg)=0$ Our right side equals: $\displaystyle\lim_{n \to \infty}(1+\frac{x}{n})^n=1+x+1*\frac{x^2}{2!}+1*\frac{x^3}{3!}+...+0+0+0+0$ $\displaystyle\lim_{n \to \infty}(1+\frac{x}{n})^n=1+x+\frac{x^2}{2!}+\frac{x^3}{3!}+...$ We got the same elements like in the Taylor series of $e^x$. Q.E.D.
2019-05-19T16:48:45
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http://bootmath.com/what-is-n-when-n0.html
What is $n!$ when $n=0$? Possible Duplicate: Prove $0! = 1$ from first principles Why does 0! = 1? If I’m right, factorial $!$ means: $$n!=1 \cdot 2 \cdot 3 \cdot 4 \cdots n$$ so: \begin{align} 5!&=1\cdot2\cdot3\cdot4\cdot5=120\\ 4!&=1\cdot2\cdot3\cdot4=24\\ 3!&=1\cdot2\cdot3=6\\ 2!&=1\cdot2=2\\ 1!&=1 \end{align} But what is $n!$ when $n=0$? It can’t be undefined and it can’t be $n!=0$, since those are illegal in known equations like: $$e^x = \sum_{n=0}^\infty \frac{x^n}{n!}$$ So what is it? Solutions Collecting From Web of "What is $n!$ when $n=0$?" $0! = 1$ is consistent with, and for reasons related to, how we define the empty product. See this entry on empty product. Empty product: The empty product of numbers is the borderline case of product, where the number of factors is zero, i.e. the set of the factors is empty. In such a “borderline” case, the empty product of numbers is equal to the multiplicative identity number, $1.$ Some of the most common examples are the following: • The zeroth power of a number x: $x^0 = 1$ • The factorial of $0: 0! = 1$ • The prime factor presentation of unity, which has no prime factors Just as ${n^0 = 1}$ for any $n$, we define, as a convention, $0!$ to be $1$. $$e^x = 1 + \frac {x} {1!} + \frac {x^2} {2!} + \frac {x^3} {3!} + … = 1 + \sum_{n=1}^\infty \frac{x^n}{n!} \tag{1}$$ But the following is a more concise definition: $$e^x = \frac {x^0} {0!} + \frac {x^1} {1!} + \frac {x^2} {2!} + \frac {x^3} {3!} + … = \sum_{n=0}^\infty \frac{x^n}{n!}\tag{2}$$ $(1)$ and $(2)$ are equal if and only if $$\;\;\displaystyle e^0 = \frac{x^0}{0!} = \frac {1}{0!} = 1 \iff 0! = 1.$$ It’s conventionally defined as $0! = 1$. This agrees with the gamma function $\Gamma(1) = (n-1)! = 1$. We know $\binom n r=\frac{n\cdot(n-1)\cdots (n-r+1)}{r!}=\frac{n!}{r! (n-r)!}$ —> the number of ways we can choose $r$ elements from $n$ elements. If $r=n,$ we can take $n$ elements from $n$ elements in $\frac{n\cdot(n-1)\cdots (n-n+1)}{n!}=1$ way. So, $1=\frac{n!}{n! (n-n)!}=\frac 1{0!}\implies 0!=1$ Similarly, if $r>n, \binom n r=\frac{n\cdot(n-1)\cdots (n-r+1)}{r!}=0\implies \frac1 {(n-r)!}=0 \implies \frac 1{s!}=0$ if $s=n-r<0$ $0!=1$. Reason 1: $(n-1)!=\dfrac{n!}n$, so $0!= \frac{1!}1$. Reason 2: $n!$ is the number of bijections of a set of cardinality $n$. The only set of cardinality $n$ is the empty set, the number of functions from the empty set to the empty set is 1. Reason 3: $n! = \int_0^{\infty} x^ne^{-x}dx$. The value for $n=0$ is 1 (and this is actually used as base of the induction proof). The number of ways to permute a set of $n$ objects is $n!$ (including the identity permutation). Your question can be reinterpreted in the following way: How many ways can one permute the elements of the empty set? Since the question can be viewed as ill-formed, one answers by conventionally defining the number to be $1$. That is, $0! = 1$. Compare this to computing the number of maps from a set of $m$ elements to a set of $n$ elements, $n^{m}$, in the case both sets are empty. Again, $0^{0} = 1$. $n!$ = $n(n-1)!$ $1!$ = $1(0)!$ $0!$ = $1$
2018-06-18T11:59:34
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https://math.stackexchange.com/questions/2611476/counterexample-for-false-o-notation-statement
# Counterexample for false O notation statement Stumbled upon some O notation counterexamples, but can't seem to figure out which one would disprove this statement: For any positive constant $c$, $f(cn)\in O(f(n))$. Any help would be appreciated! Hint: $f(cn) \in O(f(n))$ is equivalent to $\left|\frac{f(cn)}{f(n)}\right|$ being bounded. Hint: try an exponential function $f$. Let $f(n)=2^n$. Then $f(cn)=2^{cn}$ so $\frac {f(cn)}{f(n)}$ is not bounded.
2019-10-24T01:45:54
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https://math.stackexchange.com/questions/2363755/there-are-10-different-books-in-a-shelf-find-the-number-of-ways-of-selecting
There are $10$ different books in a shelf. Find the number of ways of selecting $3$ books when exactly two are consecutive. There are $10$ different books in a shelf. Find the number of ways of selecting $3$ books when exactly two are consecutive. Let us represent the problem as follows: $x_1$||$x_2$|$x_3$, where the bars are the chosen books. So $x_1+x_2+x_3=7$, where $x_1,x_3\ge 0,x_2>0$. So the number of ways are $\binom{8}{2}$ but the answer is $\binom{8}{3}$. I don't know where I am wrong. • What if the two consecutive chosen books are to the right of the other one? – aPaulT Jul 19 '17 at 11:45 You have not taken into account whether the pair of selected books appears before or after the single selected book. Thus, you must double your answer. Notice that $$2\binom{8}{2} = \binom{8}{3}$$ I would solve the problems as follows: If the consecutive books are at the ends of the shelf, then there are $10-3 = 7$ available spaces for the third book to be chosen. If the consecutive books are somewhere in the middle of the shelf, then there are $10-4 = 6$ available spaces for the third book. Since there are $2$ end spaces and $7$ spaces in the middle, the total number of ways to select the books are $6\times 7+ 7\times 2 = 56 = \binom{8}{3}$ Total number of ways $-$no consecutive $-$ all consecutive $$\binom{10}{3}-\binom{8}{3}-8=56$$
2019-05-22T04:42:22
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