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if two triangles are equal in area, they are congruent
26 de janeiro de 2021, às 3:11
They can have the same angles but have sides of different lengths. Now, the area of triangle ABC is AB*CM/2 and the area of triangle ABD is AB*DN/ 2 But CM=DN so these two areas are equal. In the simple case below, the two triangles PQR and LMN are congruent because every corresponding side has the same length, and every corresponding angle has the same measure. (Why? In this case, two triangles are congruent if two sides and one included angle in a given triangle are equal to the corresponding two sides and one included angle in another triangle. find the price at which it was sold. The two triangles have two angles congruent (equal) and the included side between those angles congruent. You can replicate the SSS Postulate using two straight objects -- uncooked spaghetti or plastic stirrers work great. (ii) If two squares have equal areas, they are congruent. #2 That is true, they must both have the same side length to have the same perimeter, therefore they will also have the same area. You now have two triangles, △SAN and △SWA. Definition: Triangles are congruent when all corresponding sides and interior angles are congruent.The triangles will have the same shape and size, but one may be a mirror image of the other. 2 triangles are congruent if they have: exactly the same three sides and SSSstands for "side, side, side" and means that we have two triangles with all three sides equal. If we can show, then, that two triangles are congruent, we will know the following: 1) Their corresponding sides are equal. An included side is the side between two angles. Two triangles are congruent if their corresponding sides are equal in length and their corresponding interior angles are equal in measure. Join now. We could also think this angle right over here. So also, sides d and e must be equal because DEF is isosceles. It can be shown that two triangles having congruent angles (equiangular triangles) are similar, that is, the corresponding sides can be proved to be proportional. I could have one triangle with sides measuring 7 in + 5 in + 8 in (perimeter = 20 in) and another triangle with sides of 6 in + 4 in + 10 in (perimeter = 20 in). Triangles are congruent when all corresponding sides & interior angles are congruent. Only if the two triangles are congruent will they have equal areas. The Angle Side Angle Postulate (ASA) says triangles are congruent if any two angles and their included side are equal in the triangles. Theorem 1 : Hypotenuse-Leg (HL) Theorem. So we have: a=d. Hence all squares are not congruent. ... maths. = True (iii) If two figures have equal areas, they are congruent. Two triangles, ABC and A′B′C′, are similar if and only if corresponding angles have the same measure: this implies that they are similar if and only if the lengths of corresponding sides are proportional. We use the symbol ≅ to show congruence. they have to be congruent to have same perimeters AND area con. …, how much loss and profit computed on his original capital. Technically speaking, that COULD almost be the end of the proof. When triangles are similar, they could be congruent. prove that they are congruent. Move to the next side (in whichever direction you want to move), which will sweep up an included angle. AAS (Angle-Angle-Side): If two pairs of angles of two triangles are equal in measurement, and a pair of corresponding non-included sides are equal in length, then the triangles are congruent. The congruence of two objects is often represented using the symbol "≅". ALL of this is based on a single concept: That the quality that we call "area" is an aspect of dimensional lengths and angles. Light1729 Light1729 If two triangles are congruent then they must overlap each other completely, so, they have exactly equal area. Yes, they are similar. And for sensible cases. Only if the two triangles are congruent will they have equal areas. Furthermore, your question is about congruent triangles and not similar triangles.l. If two triangles have the same area, they must be congruent. (d) if two sides and any angle of one triangle are equal to the corresponding sides and an angle of another triangle, then the triangles are not congruent. You can check polygons like parallelograms, squares and rectangles using these postulates. Log in. e.g. …, nswer then don't give) give me with full steps.don't spam ❌❌❌❌❌. Asked on January 17, 2020 by Premlata Pandagre In a squared sheet, draw two triangles of equal areas such that (i) the triangles are congruent. However, different squares can have sides of different lengths. Two or more triangles that have the same size and shape are called congruent triangles.. Thus, a=d. Two triangles are congruent if their corresponding sides are equal in length and their corresponding interior angles are equal in measure. Asked on December 26, 2019 by Yugandhara Jhunjhunwala. #2 That is true, they must both have the same side length to have the same perimeter, therefore they will also have the same area. All the sides of a square are of equal length. Click hereto get an answer to your question ️ If the area of two similar triangles are equal, prove that they are congruent. Triangles can be considered congruent if following conditions are satisfied. All three triangle congruence statements are generally regarded in the mathematics world as postulates, but some authorities identify them as theorems (able to be proved). (iv) If two triangles are equal in area, they are congruent. The two are not the same. one could look upside-down compared to the other but if they have the same dimensions (i.e. monojmahanta4 monojmahanta4 09.04.2020 Math Secondary School If the area of two triangles are equal. (iii) If two rectangles have equal area, they are congruent. But that does not mean that they have to be congruent. = False (iv) If two triangles are equal in area, they are congruent. That is false..one could have side lengths of 9 and 1, and the other could have side lengths of 3. See the included side between ∠C and ∠A on △CAT? You can specify conditions of storing and accessing cookies in your browser, If the area of two triangles are equal. Testing to see if triangles are congruent involves three postulates, abbreviated SAS, ASA, and SSS. An included angleis an angle formed by two given sides. (ii) If two squares have equal areas, they are congruent. Triangles can be similar or congruent. Ask your question. You can't do it. For example: (See Solving SSS Trianglesto find out more) You may have to rotate one triangle, to make a careful comparison and find corresponding parts. Join now. Why should two congruent squares have the same area? Congruent Triangles. Here, instead of picking two angles, we pick a side and its corresponding side on two triangles. Those are the three magnitudes of plane geometry: length (the sides), angle, and area. Want to see the math tutors near you? Testing to see if triangles are congruent involves three postulates. Corresponding sides and angles mean that the side on one triangle and the side on the other triangle, in the same position, match. Contrapositive of the given statement : If the areas of two traingles are not equal then the triangles are not congruent. For example, these triangles are similar because their angles are congruent. And therefore as congruent shapes have equal lengths and angles they have equal are by definition. So the fact that two triangles have the same area does not necessarily tell us that they are congruent (they could be congruent but that is often not the case) since they can have different measures on their sides or a different size and still have the same area. The triangle which maximises area for a given perimeter is the equilateral triangle with each side a third of the perimeter; the area is then p 2 /√432. (e) There is no AAA congruence criterion. Both the areas would be 9, but the figures would not be congruent. Cut a tiny bit off one, so it is not quite as long as it started out. the same shape and size, or to put it another way, the same angles and side lengths), they are congruent. So I guess the answer is not always. (f) Two circles having same circumference are congruent. Compare them to the corresponding angles on △BUG. Introducing a diagonal into any of those shapes creates two triangles. Definition: Triangles are congruent when all corresponding sides and interior angles are congruent.The triangles will have the same shape and size, but one may be a mirror image of the other. If the area of two similar triangles are equal, prove that they are congruent. The given statement -. So let's just start That triangle BCD is congruent Let’s use congruent triangles first because it requires less additional lines. -If two squares have equal perimeters, then they have equal ...” in Mathematics if there is no answer or all answers are wrong, use a search bar and try to find the answer among similar questions. The SAS Postulate says that triangles are congruent if any pair of corresponding sides and their included angle are congruent. Congruent triangles will have completely matching angles and sides. The meaning of congruence in Maths is when two figures are similar to each other based on their shape and size. If the areas of two similar triangles are equal, prove that they are congruent. Put them together. Pick any side of △JOB below. Only if the two triangles are congruent will they have equal areas. Prove that if in two triangles,two angles and the included side of one triangle are equal to two angles and the included side of the other triangle,then two triangles are congruent. the cost price of a flower vase is rs 120 . The diagonals of a rhombus divide it into four triangles of equal area. Converse of the above statement : If the areas of the two triangles are equal, then the triangles are congruent. Congruence is our first way of knowing that magnitudes of the same kind are equal. (iv) If two triangles are equal in area, they are congruent. For example, both of these triangles are isosceles, since they have two equal sides and angles. So go ahead; look at either ∠C and ∠T or ∠A and ∠T on △CAT. So once you realize that three lengths can only make one triangle, you can see that two triangles with their three sides corresponding to each other are identical, or congruent. Congruent triangles are triangles which are identical, aside from orientation. By applying the Side Angle Side Postulate (SAS), you can also be sure your two triangles are congruent. Now you have three sides of a triangle. Divide a square sheet in two triangles of equal area so that they are congruent. If the corresponding angles of two triangles are equal, then they are always congruent. … Using any postulate, you will find that the two created triangles are always congruent. 0 votes . Congruent Triangles. If 2 sides and the included angle of one triangle are equal to the corresponding sides and angle of another triangle, the triangles are congruent. Hence, the congruence of triangles can be evaluated by knowing only three values out of six. If the hypotenuse and one leg of a right triangle are equal to the hypotenuse and one leg of another right triangle, then the two right triangles are congruent. What is the solution set of y=3x+16 and 2y=x+72? So, for equal area, all sides are equal. asked Jul 7 in Congruence of Triangles by Rani01 (52.4k points) congruent triangles; class-7; 0 votes. Two triangles are congruent if they have the same shape and size, but their position or orientation may vary. You may think we rigged this, because we forced you to look at particular angles. If you're behind a web filter, please make sure that the domains *.kastatic.org and *.kasandbox.org are unblocked. Think: Two congruent triangles have the same area. ASA (Angle-Side-Angle): If two pairs of angles of two triangles are equal in measurement, and the included sides are equal in length, then the triangles are congruent. Side-Angle-Sideis a rule used to prove whether a given set of triangles are congruent. (if you don't know the a If the areas of two similar triangles are equal, prove that they are congruent. Geometricians prefer more elegant ways to prove congruence. But just to be overly careful, let's compute a/d. $\endgroup$ – Moti Feb 7 '15 at 16:42 add a comment | Your Answer In the sketch below, we have △CAT and △BUG. State true or false and justify your answer. In most systems of axioms, the three criteria – SAS, SSS and ASA – are established as theorems. (iii) If two figures have equal areas, they are congruent. Since two triangles are similar therefore the ratio of the area is equal to the square of the ratio of its corresponding side a r e a ∆ A B C a r e a ∆ D E F = B C E F 2 = A B D E 2 = A C D F 2 B C E F 2 = A B D E 2 = A C D F 2 = 1 N o w , t a k i n g a n y o n e c a s e 1 = B C E F 2 1 = B C E F E F = B C (ii) If two squares have equal areas, they are congruent. In the simple case below, the two triangles PQR and LMN are congruent because every corresponding side has the same length, and every corresponding angle has the same measure. prove that they are congruent. ... triangles are equal, prove that they are congruent. Those are the three magnitudes of plane geometry: length (the sides), angle, and area. You can compare those three triangle parts to the corresponding parts of △SAN: After working your way through this lesson and giving it some thought, you now are able to recall and apply three triangle congruence postulates, the Side Angle Side Congruence Postulate, Angle Side Angle Congruence Postulate, and the Side Side Side Congruence Postulate. You can think you are clever and switch two sides around, but then all you have is a reflection (a mirror image) of the original. Congruent triangles are triangles having corresponding sides and angles to be equal. https://www.mathsisfun.com › geometry › triangles-congruent.html If the areas of two similar triangles are equal, prove that they are congruent. An included angleis an angle formed by two given sides. But, it is not necessary that triangles having equal area will be congruent to one another they may or may not be congruent. Conditional Statements and Their Converse. For eg., [ example for non-congruent triangles having equal areas] (1) ABC has base BC = 6 cm and height AD = 4 cm, other sides AB =12 cm, CA = 17 cm. Now you will be able to easily solve problems on congruent triangles definition, congruent triangles symbol, congruent triangles Class 8, congruent triangles geometry, congruent t If we can show, then, that two triangles are congruent, we will know the following: 1) Their corresponding sides are equal. A, B and C are three sets, n(A) = 14, n(C) = 13, n(A U BUC) = 28, Local and online. This forces the remaining angle on our △CAT to be: This is because interior angles of triangles add to 180°. ASA (angle, side, angle) ASA stands for “ angle, side, angle ” and means that we have 2 triangles where we know 2 angles and the included side are equal. Similar triangles will have congruent angles but sides of different lengths. are both 90 degrees. Guess what? Congruence is denoted by the symbol ≅. If triangle RST is congruent to triangle WXY and the area of triangle WXY is 20 square inches, then the area of triangle RST is 20 in.². The four triangles are congruent with each other regardless whether they are rotated or flipped. Find n(ABC).9. The given statement - "If two triangles are congruent, then their areas are equal." Learn faster with a math tutor. …, n(A) = n(B), nn Byn(B, C) = 4 and n(An C) = 3. This is the only postulate that does not deal with angles. AAS is equivalent to an ASA condition, by the fact that if any two angles are … If you have two similar triangles, and one pair of corresponding sides are equal, then your two triangles are congruent. Log in. A postulate is a statement presented mathematically that is assumed to be true. Two right triangles can be considered to be congruent, if they satisfy one of the following theorems. To be congruent two triangles must be the same shape and size. asked by priyanka. Congruent triangles sharing a common side. Let's take a look at the three postulates abbreviated ASA, SAS, and SSS. You will see that all the angles and all the sides are congruent in the two triangles, no matter which ones you pick to compare. After you look over this lesson, read the instructions, and take in the video, you will be able to: Get better grades with tutoring from top-rated private tutors. If the hypotenuse and one leg of a right triangle are equal to the hypotenuse and one leg of another right triangle, then the two right triangles are congruent. You can only assemble your triangle in one way, no matter what you do. FALSE. Triangle Congruence Theorems (SSS, SAS, ASA), Congruency of Right Triangles (LA & LL Theorems), Perpendicular Bisector (Definition & Construction), How to Find the Area of a Regular Polygon, Do not worry if some texts call them postulates and some mathematicians call the theorems. "If two triangles are congruent, then their areas are equal." b=e. Pairs - The classic pairs game with simple congruent shapes. So now you have a side SA, an included angle ∠WSA, and a side SW of △SWA. Note that if two angles of one are equal to two angles of the other triangle, the tird angles of the two triangles too will be equal. • SSS (Side Side Side) if all three corresponding sides are equal in length. 3) They are equal areas. it the shopkeeper sells it at a loss of 10% . TRUE. The legs of each of these isosceles triangles could have any lengths as long as they are equal, but the legs of these two triangles need not be the same. CPCT Rules in Maths. Yes, if two triangles have two congruent angles and two congruent sides then the triangles are guaranteed to be congruent. It is true that all congruent triangles have equal area. Comparing one triangle with another for congruence, they use three postulates. https://www.mathsisfun.com/geometry/triangles-congruent.html Remember that the included angle must be formed by the given two sides for the triangles to be congruent. Illustration of SAS rule: Suppose you have parallelogram SWAN and add diagonal SA. 2) Their corresponding angles are equal. However, they do not have to be congruent in order to be similar. asked Sep 21, 2018 in Class IX Maths by navnit40 ( -4,939 points) That is false..one could have side lengths of 9 and 1, and the other could have side lengths of 3. This will happen when the area is half the multiplication of the sides, or maximum area for such two sides. -If two rectangles have equal areas, then they are congruent. The postulate says you can pick any two angles and their included side. You can now determine if any two triangles are congruent! Find a tutor locally or online. Side-Angle-Sideis a rule used to prove whether a given set of triangles are congruent. Notice that ∠C on △CAT is congruent to ∠B on △BUG, and ∠A on △CAT is congruent to ∠U on △BUG. (g) If two triangles are equal in area, they are congruent. TRUE. Corresponding sides and angles mean that the side on one triangle and the side on the other triangle, in the same position, match. For two triangles to be congruent, the corresponding angles and the sides are to be equal. Perhaps the easiest of the three postulates, Side Side Side Postulate (SSS) says triangles are congruent if three sides of one triangle are congruent to the corresponding sides of the other triangle. Are they congruent? 1. You can only make one triangle (or its reflection) with given sides and angles. Two right triangles can be considered to be congruent, if they satisfy one of the following theorems. Write the following statements in words, farhan sarted his bussiness by investing 75000 in the first year he made a profit of 20% he invested the capital in new business and made loss of 12% prove that they are congruent. 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https://www.physicsforums.com/threads/quick-question-about-integral-of-1-x.300986/ | # Quick question about integral of (1/x)
1. ### 2^Oscar
45
I know that the integral of 1/x is ln(x)
but if you follow the conventional method of integrating you raise the power and multiply by the new power, and you get an answer that is obviously wrong
Would someone mind explaining to me why you cannot integrate this function by normal means?
Thanks,
Oscar
2. ### jdougherty
25
One way to think of it is the following: If we could find the antiderivative by "normal" means, then we would get
$\int \frac{1}{x}~dx = x^{0} = \textrm{constant}$
which is clearly not true from looking at the graph of 1/x
Another way to look at it: the reason we can do the reverse power rule usually for two reasons: 1) the fundamental theorem of calculus, and 2) the power rule. In other words, it usually works because we know if we find the function F(x) such that $x^{n}$ is the derivative, then we know the integral is equal to F(x). But the power rule only works when n is not zero, since D[$x^{0}$] = D[1] = 0.
3. ### lurflurf
2,325
This is silly.
You are maybe thinking of
Integral[x^a]=x^(a+1)/(a+1) (* a!=-1)
This cannot hold when a=-1 because it would require division by zero.
We can however use limits to extend the result.
Integral[x^a]=lim_{a->a} x^(a+1)/(a+1)
in particular
Integral[x^-1]=lim_{a->-1} x^(a+1)/(a+1)=log(x)
Last edited: Mar 20, 2009
4. ### lanedance
3,307
hmmm.. do you mean
$$\int x^a = ^{lim}_{b\rightarrow a} \frac{x^{b+1}}{b+1}$$
then
$$\int \frac{1}{x}= ^{lim}_{b\rightarrow -1} \frac{x^{b+1}}{b+1} =ln(x)$$
5. ### HallsofIvy
40,310
Staff Emeritus
It is not clear what you are asking. The obvious answer is what lurflurf said. You cannot use the formula
$$\int x^n dx= \frac{1}{n+1} x^{n+1}+ C$$
when n= -1 because then you would be dividing by 0. I'm not sure why you think of that as "normal means".
6. ### confinement
192
Think of it this way, the antiderivative of 1/x is the function whose inverse is exactly equal to its own derivative. Let y(x) be the antiderivate of 1/x. Then we have:
$$\frac{dy}{dx} = \frac{1}{x}$$
Inverting the Liebniz notation in the way that he intended yields:
$$\frac{dx}{dy} = x$$
The last equation says that x' = x, i.e. the function x(y) is equal to its own derivative. This means that x(y) cannot be a polynomial or rational function, since all of those functions change when you differentiate them. It is this perfect property, that the antiderivative of 1/x is the inverse function of the function who is exactly equal to its own derivative, that puts in a class of its own, a special case.
7. ### 2^Oscar
45
Hi,
Sorry for being unclear.
My question was supposed to ask why you cannot do the following:
$$\int x^n dx= \frac{1}{n+1} x^{n+1}+ C$$
for the special case of x-1, I realise that it cannot be done from looking at the graph of y=$$1/x$$ and because the fraction requiring division by 0 cannot be defined. my query was as to why you could not do this for x-1 (which confinement answered).
Once again appologies for the lack of clarity in my question, all your responses have cleared some of the clouds from my head :)
thanks very much,
Oscar
8. ### lurflurf
2,325
Yes, though lim a->a is valid it is confusing to some, I also typoed a into x. Also omiting dx is valid, but confusing to some. Most people denote their favorite base as log, what then is you favorite base if not e?
9. ### lanedance
3,307
yeah its got to be e, though i like the ln just to avoid any confusion
Last edited: Mar 21, 2009
10. ### nikkor180
5
Greetings:
d/dx [ln(x)] = lim(h-->0) [(ln(x+h) - ln(x)) / h]
= lim(h-->0) [(1/h)*ln[(x+h)/x]
= (1/x) lim(h-->0) [(x/h) ln(1+h/x)]
= (1/x) lim(h-->0) [ln(1+h/x)^(x/h)]
= (1/x) ln [lim(h-->0) [(1+(h/x))^(x/h)]] {from lim[f(g(x))] = f(lim[g(x)])}
= (1/x) ln e
= 1/x.
Since d/dx [ln(x)] = 1/x, it follows that intgrl[1/x] = ln x.
Regards,
Rich B.
[email protected]
11. ### arildno
12,015
You might prefer the following argument for why the anti-derivative should be ln(x) plus some constant C:
1. We define the arbitrary (non-natural) power of some (positive) variable as follows:
$$x^{a}=e^{a\ln(x)}$$
2. In terms of natural powers, we define the exponential function itself as follows:
$$e^{y}=1+\sum_{n=1}^{\infty}\frac{y^{n}}{n!}$$
3. Let us utilize the following definite integral to prove our point:
We have:
$$\int_{a}^{b}x^{-1+\epsilon}dx=\frac{1}{\epsilon}(b^{\epsilon}-a^{\epsilon})$$,
where a,b,$\epsilon$ are all considered greater than 0.
4. Now, let us utilize definitions from 1, 2 upon the expression in 3. We gain:
$$\frac{1}{\epsilon}(b^{\epsilon}-a^{\epsilon})=\sum_{n=1}^{\infty}\frac{\epsilon^{n-1}((\ln(b))^{n}-(\ln(a))^{n}}{n!})$$
5. Thus, in the limit $\epsilon\to{0}$, all but the first double-term of this sum goes to zero, and we retain:
$$\lim_{\epsilon\to{0}}\int_{a}^{b}x^{-1+\epsilon}dx=ln(b)-ln(a)$$
This is in accordance with what we should have!
6. Note that this argument shows that the definite integral of some power of x can be defined as a continuous function of the power variable, even in the "special case" where the power variable has the value -1.
Last edited: Mar 21, 2009
12. ### sbcdave
10
I'm new to calculus, but my impression is that an integral should be a function that represents the area under the curve of the function being integrated, which f(x)=ln(x) does not do for f(x)=1/x
In the graph of 1/x you can see that from 0 toward infinity the area under the curve would come on quickly and almost stop increasing.
Can anyone shed light on what I'm missing.
I apologize if this is elementary. I took college algebra 10 years ago, recently took a precalc class at a community college and just finished reading a text book called Brief Calculus.
Love math and physics by the way and am glad I found this site.
13. ### checkitagain
99
To 2^Oscar and the rest of this thread's users:
--------> $\int{\dfrac{1}{x}}$dx$\ = \ \ln|x| \ + \ C$
Other than missing the arbitrary constant, "C," there must
be absolute value bars around the "x."
$*** Edit ***$
$\text{Good catch by micromass concerning the missing dx term.}$
Last edited: Dec 21, 2011
14. ### lavinia
1,989
While the answers already given are correct here is a way of looking at it that has always mystified me.
1/x has the remarkable property that the area under it from 1 to some number is the sum of the areas from 1 to any two numbers whose product equals the number. So if x = yz then the area from 1 to x is the sum of the areas from 1 to y and 1 to z. It is easy to show this using the integral definition of area though I would love to see an elementary proof not using calculus.
It is not hard to convince yourself that a rational function (quotient of two polynomials) can not have this property. So the logarithm is something new. It's derivative is a rational function but it is not. It is like a portal from the rational to the transcendental.
A function that compares two laws of addition and/or multiplication is called a homomorphism. The logarithm is a homomorphism from the positive numbers under multiplication to all of the numbers under addition. The exponential function is the reverse homomorphism taking the numbers under addition to the positive numbers under multiplication.
Last edited: Dec 21, 2011
15. ### checkitagain
99
$\text{Not those answers as they concern where the}$
$\text{absolute value symbol is missing, not to mention the "+ C."}$
Please see the post above yours regarding this.
16. ### micromass
18,553
Staff Emeritus
You forgot the dx in your post above. So your answer is also incorrect.
17. ### dextercioby
12,304
Somehow (to me at least) related to lavinia's post is the interesting fact discovered by Euler
$$\mbox{for x very large} ~ \ln x \simeq \frac{1}{1} + \frac{1}{2} +\frac{1}{3} +...+\frac{1}{x}$$
a result which can be recast rigorously in terms of limits.
18. ### lanedance
3,307
Hey sbcdave -welcome to PF!
you're generally best off posting question as new threads, gets less confusing that way, particularly for old threads.
The reminann integral (there are other more complex defintions) is interpreted as the area under of a well behaved function
With this is mind and considering it as a definite integral (over a given integral where the function is well behaved
$$\int_a^bdx \frac{1}{x} = ln(x)|_a^bdx =ln(b)-ln(a)$$
This is equivalent to the area between the function and the horizontal axis. Notice ln(x) is a a "monotonically" increasing function (ln(b)>ln(a) for all 0<a<b), so it always gives a positive area.
I don't totally undertstand your question..
But, at x = 0, f(x) = 1/x is not well defined, so to calculate the area you must use limits.
Similarly to deal with taking the integral to infinity, you need to use a limiting process.
Things get subtle when you consder lmiting process, so its important to be rigrorous. First lest split the integral into 2 portions, which we can do when the function is well behaved:
$$\int_a^bdx f(x) = \int_a^cdx f(x) + \int_c^bdx f(x)$$
In this case lets choose c=1 and the integral becomes
$$\lim_{a \to 0^+} \lim_{b \to +\infty} \int_a^b dx \frac{1}{x} = \lim_{a \to 0^+} \int_a^1 dx \frac{1}{x} + \lim_{b \to +\infty}\int_1^bdx \frac{1}{x} = \lim_{a \to 0^+}(ln(1)-ln(a)) + \lim_{b \to +\infty}(ln(b)-ln(1)) = \lim_{a \to 0^+}{-ln(a)} + \lim_{b \to +\infty}ln(b)$$
As both these diverge (become infinite) we actually find there is inifinite area below 1/x in both the intervals (0,1) and (1,inf). In fact in some repsects there is saimialr amount of infinte area in each.
Though in both cases the curve compresses against the axis, it doesn't do so quickly enough to result in a finite area - things can get wierd with lmiting processes
Hopefully this helped answer your question and didn't confuse things more!
can be a nice way to help, learn and keep up with some math skills
Mod note: fixed LaTeX
Last edited by a moderator: Dec 21, 2011
19. ### sbcdave
10
Thanks
Roger
But the natural log of anything less than 1 is negative, which seems to say that the area under the curve from 0 to anything less than x=1 on the graph of f(x)=1/x should be a negative area. The graph of f(x)=1/x from x=0 to 1 is above the x-axis and as you say in the next quote is infinite.
Example of what I'm struggling with: $$\int_{0.1}^1dx \frac{1}{x}$$
Should be a definable positive integer. Approximately equal to 9(0.1) increments with heights 1/0.9 + 1/0.8 + ... 1/0.1 (left hand method) or approximately equal to 2.829. However, if you do ln(1)-ln(0.1) you get approximately 2.3 (I expected a negative number and need to think about this a little more, there's clearly a flaw in my logic...haha) Sorry. Posting anyways so that you might see where I was coming from.
20. ### sbcdave
10
I was looking at a graph of ln(x) and 1/x together on my ti-83 and assumed that because the ln(x) went negative with x<1 that that meant the area was negative, now I see that because the -y value grows as x gets smaller that your subtracting a larger negative from a negative and end up with a positive.
Now I'm a little confused why the left hand method was larger than the ln(x)-ln(x) method instead of smaller though if you have any ideas on that?
Thanks again for humoring me once already I'll understand if the thread dies here...lol
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Have something to add? | 2015-04-01T01:08:50 | {
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https://math.stackexchange.com/questions/1156899/proving-a-open-in-v-subseteq-m-metric-space-iff-a-c-cap-v-certain-open | # Proving '$A$ open in $V\subseteq M$ (metric space) iff $A=C\cap V$ (certain open $C$ in $M$)'
I want to prove the following:
Let $(M,d)$ be a metric space. Let $A\subseteq V\subseteq M$.
1) $A$ is open in $V \Leftrightarrow A = C\cap V$ (for a certain open $C$ in $M$)
2) $A$ is closed in $V \Leftrightarrow A = C\cap V$ (for a certain closed $C$ in $M$)
### Questions:
• Could someone check the proof?
• 'for a certain open $C$ in $\color{Blue}{M}$.'
Would this proof also work for a more specific choice of $C$? Like for a certain open $C$ in $\color{blue}{V}$. I don't really see the added value of choosing $M$ over $V$.
• Could some give me some pointers on how to prove $2, \Rightarrow$?
### Proof 1)
$\Leftarrow$: Choose $a\in A$.
$$\begin{array}{rl} & a \in A = C\cap V\\ \Rightarrow & a \in C\\ \Rightarrow & (\exists r > 0)(B_M(a,r)\subseteq C)\\ \Rightarrow & (\exists r > 0)(B_M(a,r)\cap V \subseteq C\cap V)\\ \Rightarrow & (\exists r> 0) (B_V(a,r)\subseteq A \end{array}$$
$\Rightarrow$: Choose $a\in A$.
$$\begin{array}{rl} \Rightarrow & (\exists r_a >0)(B_V(a,r_a) \subseteq A) \end{array}$$
Consider all $a\in A$ then:
$$\begin{array}{rl} & A = \bigcup_{a\in A} B_V(a,r_a)\\ \Rightarrow & A = \bigcup_{a\in A} \left[ V\cap B_M(a,r_a)\right]\\ \Rightarrow & A = V\cap\left[ \bigcup_{a\in A} B_M(a,r_a)\right] \end{array}$$
Let $$\left[ \bigcup_{a\in A} B_M(a,r_a)\right] = C$$ which is open as a union of open sets.
### Proof 2)
$\Leftarrow$:
$$\begin{array}{rrl} & V\setminus A &= V\setminus(C\cap V)\\ \Rightarrow & & = (V\setminus C)\cup (V\setminus V)\\ \Rightarrow && = V\setminus C \end{array}$$
Since $C$ is closed then $V\setminus C$ is open and so is $V\setminus A$. Then $A$ is closed in $V$.
$\Rightarrow$: How?
So far everything seems right. For the last proof, let $A$ closed in $V$. Then $V\backslash A$ is open in $V$. By the first statement (which you already proved) there exists an open $B$ in $M$ such that $V\backslash A=B\cap V$. Now $A=(M\backslash B)\cap V$. And we know that $C= M\backslash B$ is closed in $M$ because $B$ is open in $M$. | 2019-09-23T05:05:28 | {
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https://math.stackexchange.com/questions/2221033/is-0-an-imaginary-number | # Is $0$ an imaginary number?
My question is due to an edit to the Wikipedia article: Imaginary number.
The funny thing is, I couldn't find (in three of my old textbooks) a clear definition of an "imaginary number". (Though they were pretty good at defining "imaginary component", etc.)
I understand that the number zero lies on both the real and imaginary axes.
But is $\it 0$ both a real number and an imaginary number?
We know certainly, that there are complex numbers that are neither purely real, nor purely imaginary. But I've always previously considered, that a purely imaginary number had to have a square that is a real and negative number (not just non-positive).
Clearly we can (re)define a real number as a complex number with an imaginary component that is zero (meaning that $0$ is a real number), but if one were to define an imaginary number as a complex number with real component zero, then that would also include $0$ among the pure imaginaries.
What is the complete and formal definition of an "imaginary number" (outside of the Wikipedia reference or anything derived from it)?
• See Complex number : "A complex number whose real part is zero is said to be purely imaginary, whereas a complex number whose imaginary part is zero is a real number." Apr 6 '17 at 15:51
• Here's what wolfram says Apr 6 '17 at 15:52
• I do not think this question should be down voted. It is well edited and clearly there was decent thought put into it.
– user416426
Apr 6 '17 at 15:58
• The downvotes are sad. The premise might seem silly, but the question is well-written and clearly thought-out. I like it. Apr 6 '17 at 16:02
• And why not? Mathematics is full of similar cases. For example, the zero function is the unique function that is both even and odd.
– MJD
Apr 6 '17 at 16:02
The Wikipedia article cites a textbook that manages to confuse the issue further:
Purely imaginary (complex) number : A complex number $z = x + iy$ is called a purely imaginary number iff $x=0$ i.e. $R(z) = 0$.
Imaginary number : A complex number $z = x + iy$ is said to be an imaginary number if and only if $y \ne 0$ i.e., $I(z) \ne 0$.
This is a slightly different usage of the word "imaginary", meaning "non-real": among the complex numbers, those that aren't real we call imaginary, and a further subset of those (with real part $0$) are purely imaginary. Except that by this definition, $0$ is clearly purely imaginary but not imaginary!
Anyway, anybody can write a textbook, so I think that the real test is this: does $0$ have the properties we want a (purely) imaginary number to have?
I can't (and MSE can't) think of any useful properties of purely imaginary complex numbers $z$ apart from the characterization that $|e^{z}| = 1$. But $0$ clearly has this property, so we should consider it purely imaginary.
(On the other hand, $0$ has all of the properties a real number should have, being real; so it makes some amount of sense to also say that it's purely imaginary but not imaginary at the same time.)
I don't think there is a
complete and formal definition of "imaginary number"
It's a useful term sometimes. It's an author's responsibility to make clear what he or she means in any particular context where precision matters. If $0$ should count, or not, then the text must say so.
Your question shows clearly that you understand the structure of the complex numbers, so you should be able to make sense of any passage you encounter.
A complex number z=a+ib where a and b are real numbers is called : 1- purely real , if b=0 ; e.g.- 56,78 ; 2- purely imaginary, if a=0 ,e.g.- 2i, (5/2)i ; 3- imaginary,if b≠ 0 ,e.g.- 2+3i,1-i,5i ; 0 is purely imaginary and purely real but not imaginary. | 2022-01-20T15:39:27 | {
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http://math.stackexchange.com/questions/205778/prove-that-a-language-b-is-regular | # Prove that a language B is regular
here is the question I'm dealing with:
Let B = {$1^{k}$y|y $\in$ {0,1}* and y contains at least k 1s, for k $\geq$ 1}.
Show that B is a regular language.
Can I use the pumping lemma to show that this is a regular language? My intuition tells me that I can't, since I would have to find every string s that satisfies the pumping lemma.
Anyways, your suggestions are appreciated. Thank you in advance.
-
The pumping lemma for regular languages is only a tool for showing that a language is not regular; it cannot be used to show that a language is regular. The two most straightforward ways to demonstrate that a language is regular are (1) to write a regular grammar that generates it, and (2) to design a finite state automaton that recognizes it.
In this case, though, it pays to begin by taking a close look at just what words are in $B$. Consider a word $11x$, where $x\in\{0,1\}^*$: since we can set $k=1$, such a word is automatically in $B$, since $1x$ certainly contains at least one $1$. A word that begins with $1$ is not in $B$ if and only if it has no other $1$’s; and a word that begins with $0$ is not in $B$. Thus,
$$B=\{1x\in\{0,1\}^*:x\in\{0,1\}^*\text{ and }x\text{ has at least one }1\}\;,$$
the language corresponding to the regular expression $10^*1(0\lor 1)^*$. It’s not at all hard to design a regular grammar that generates this language, or a finite state machine that recognizes it.
-
I just built the DFA and thus proved that B is regular. Thanks so much for confirming my initial assumption. Thanks a lot! – Dave Oct 2 '12 at 3:22
@Dave: You’re welcome! – Brian M. Scott Oct 2 '12 at 3:24
wait a second.. what about string 1101. this is not accepted by B (take k=2). New B would accept it. Am I wrong? (I would be wrong if you could always pick k=1) – Dave Oct 2 '12 at 3:50
@Dave: The way the description of $B$ is worded allows you to choose $k$. Had it said that there is some $k\ge 1$ such that $B=\{1^ky:\text{ etc. }\}$, that would have been a different story, though not much harder. If it meant that you had to use the maximum $k$ allowed by the initial string of $1$’s, it should have said so; that would be very different. – Brian M. Scott Oct 2 '12 at 3:52
@TheNotMe: That’s even better. You’re welcome! – Brian M. Scott Nov 30 '13 at 17:25
show 4 more comments
You can never use (directly) the pumping lemma to show a langauge is regular since there are non-regular langauges that satisfy all the lemma properties.
In your case the langauge is very simple and you can build a finite state automata to prove regularity
- | 2014-04-20T22:19:43 | {
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https://www.physicsforums.com/threads/change-in-scales-causing-change-in-gradient.947892/ | # Change in scales causing change in gradient?
## Main Question or Discussion Point
Hello I plotted a graph of velocity against time only to realise that I needed more space on the X axis, so I changed my scale from instead of 2cm= 10 seconds to 2cm =20 seconds.
My Y axis remained constant with a scale of 2cm= 10 m/s
However, the gradient I got from the first scale was 1.15m/s2 and this even coincides with the data given.
But with my new readjusted scale to fit the page I'm getting a gradient of 2.4 m/s2.
How can this be? the graph is a straight line so the acceleration is constant according to the data given, am I plotting this wrong??
DrClaude
Mentor
However, the gradient I got from the first scale was 1.15m/s2 and this even coincides with the data given.
How are you calculating that gradient?
Dr Claude, I am using the triangle method and then dividing the vertical units by the horizontal units, using two poitns that are NOT given in the data but instead two points from the line of best fit I drew, which is standard practice for this level
jedishrfu
Mentor
It's somewhat hard to determine what you did wrong without seeing any work or a picture of your chart.
If you've computed the slope using algebraic means you should get the same value ie $(y1-y0)/(x1-x0)$ for your acceleration (velocity/seconds ie $m/s^2$)
If instead you are using the grids on your graph paper then of course you'll get something different unless you adjust for the x length of the grid element.
As an example, if you computed it using $(y1-y0)$ in centimeters over $(x1-x0)$ in centimeters then that is wrong.
Whereas if you computed it $( (y1-y0).in.cm) * (10m/s.per.cm) )$ over $( ((x1-x0).in.cm) * (20 seconds.per.cm))$
DrClaude
Mentor
Dr Claude, I am using the triangle method and then dividing the vertical units by the horizontal units
Then my guess is that you are not calculating the units correctly.
Can you post the plots you did? It may help in figuring where you are going wrong.
Ok here are my images with the scales represented above and the data presented to plot orginally...
However I realized what I did wrong... I simply forgot to convert thehorizontal units in my new graph the second picture, to suit...it was just one tiny mistake.... how DO I avoid these? The test I'm studying for specifically wants me to do the triangle method, bigger than half the line, etc.
BUt theres not much time and I somtimes I just make silly errors, now I feel dumb sorry for this silly post
#### Attachments
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jedishrfu
Mentor
This is NOT a silly post. We all learn from our mistakes.
The most humbling thing is that we learn we make mistakes. The second thing is in reconstructing our reasoning we might determine why we made the mistake.
To be fair, you caught your mistake and that is a very hopeful sign that you're really grasping the subject. This mistake will stick with you forever and I'm sure you won't make this same mistake twice. (fingers crossed). You can even pass it on to your kids with some sagely advice.
Our mistakes teach us so much and eventually we get to teach them to others.
scottdave | 2020-04-02T10:39:47 | {
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http://spmd.safoo.de/geometric-series-test-proof.html | ## Geometric Series Test Proof
Theorem: (Geometric Series Test) If jrj<1, the geometric series P 1 n=0 ar n, where a6= 0, converges with sum a 1 r: If jrj 1, the series diverges. Let s n = be the n th partial sum. In an ecological study of the feeding behavior of birds, the number of hos between flights was counted for several birds. Now, as we have done all the work with the simple arithmetic geometric series, all that remains is to substitute our formula, (Noting that here, the number of terms is n-1). The word "countable" means that you can label. By inspection, it can be difficult to see whether a series will converge or not. Example 13. Definition. The geometric series and the telescoping series make their appearance in this chapter. (Alternating series test) Consider the series. These are both geometric series, so I can sum them using the formula for geometric series: X Series Test says that the series converges. P), then b is the arithmetic mean of a and c. Water authorities identify a threshold geometric mean where beaches or shellfish beds must be closed. The proof by induction method use to proof geometric series can be applied to other progression as well by doing the same step:. Test results for water quality (specifically, fecal coliform bacteria concentrations) are sometimes reported as geometric means. The geometric series converges and has a sum of if. If the limit of a[n] is not zero, or does not exist, then the sum diverges. 4 Ratio test The geometric series leads to a useful test for convergence of the general series X1 n=0 a n= a 0 + a 1 + a 2 + (12) We can make sense of this series again as the limit of the partial sums S n = a 0 + a 1 + + a n as n!1. The Egyptians used this method of finite geometric series mainly to "solve problems dealing with areas of fields and volumes of granaries" but used it for many other uses too, including the pyramids and math problems similar to those one might find on a STAAR test today (see D1, and F1). If we can –nd a function f(x) such that. 1)25 is a geometric progression with a geometric ratio of (1. Now if 1 k=N+2 a k converged, this implies a k!0. usually shown using the integral test. The first proof in Algebra 2! Students learn to derive the formula for the sum of the first n terms of a finite geometric sequence. 6Theorem (Alternating series test): (i) Let be an alternating series such that (ii) Then is convergent. The Mathematics Vision Project: Scott Hendrickson, Joleigh Honey, Barbara Kuehl, Travis Lemon, Janet Sutorius. The nth-term test for divergence. Any such series can be written as ± P (−1)kak with ak ≥ 0 for all k ∈ N. Finally, we give a test which helps us to analyze convergence of an alternating series. Finding Geometry Help is Easy. You can see that this is reasonable by dividing 1 by , or using the the formula for the sum of a geometric series with ratio. Buy ColorBird Geometric Series Tablecloth Diamond Pattern Cotton Linen Dust-Proof Table Cover for Kitchen Dinning Tabletop Linen Decor (Round, 60 Inch, Yellow): Tablecloths - Amazon. Hence, is absolutely convergent, and thus is itself convergent. 1 nm in a pressurized stirred cell (detailed in fig. This task also provides practice in writing and using formulas for arithmetic sequences. Finally, we give a test which helps us to analyze convergence of an alternating series. In partnership with the. In fact, our proof is an extension of the nice result given by Cohen and Knight [2]. These functions would be difficult if not impossible to write down analytically, but there is software to find conformal maps numerically. A SHORT(ER) PROOF OF THE DIVERGENCE OF THE HARMONIC SERIES LEO GOLDMAKHER It is a classical fact that the harmonic series 1+ 1 2 + 1 3 + 1 4 + diverges. Property 1: If |r| < 1 then the geometric series converges to. Free math lessons and math homework help from basic math to algebra, geometry and beyond. 1),…, 500(1. You can then "flex" the 3-D hexaflexagon, exposing each of four six-sided faces, one at a time. There are several questions Consider the geometric series: (Sum from k=0 to infinity) of ar^k and consider the repeating decimal. The circle has radius 25 centimetres and the angle of the sector is 280° a) Find the area of the sector of the circle. If L>1 then X a n diverges. The idea is to compare the series witha geometric series witha ratio slightly larger than ½. Determine the number of terms n in each geometric series. This looks to me like an expression of the "ratio test" for convergence of a series. Grappling with the geometric series, geometry formulas or geometric sequence? Our tutors can help. I we see from the graph that because the values of b n are decreasing, the. In fact this series diverges quite slowly. To make a 3-D hexaflexagon, print out the template, cut and fold carefully, then tape (or glue) into shape. If a n diverges, then b n diverges. Great! Think it might be an arithmetic or geometric sequence? If the sequence has a common difference, it's arithmetic. Radius of Convergence. Now, as we have done all the work with the simple arithmetic geometric series, all that remains is to substitute our formula, (Noting that here, the number of terms is n-1). We also consider two specific. However, our rules of probability allow us to also study random variables that have a countable [but possibly infinite] number of possible values. Tests for convergence: Ratio Test, Root Test and Raabe's Test. So, the sum of the series, which is the limit of. If 1 then 1 for all so does not define a null sequence and the series diverges by the null sequence test. Geometric Proof Video 366. In mathematics, a geometric series is a series with a constant ratio between successive terms. Direct Comparison Test If 0 <= a n <= b n for all n greater than some positive integer N, then the following rules apply: If b n converges, then a n converges. NOTES ON INFINITE SEQUENCES AND SERIES MIGUEL A. Alternating series Theorem (Leibniz’s test) If the sequence {a n} satisfies: 0 < a n, and a n+1 6 a n, and a n → 0, then the alternating series P ∞ n=1 (−1) n+1a n converges. This quiz/worksheet combination will test your understanding of the formula to find the sum of the infinite geometric series by providing you with example problems. Suppose the interest rate is loo%, so i = 1. For example, if , Results on geometric series show that the two expressions are equal. The next proof is unique among all known proofs of the infinitude of the set of primes. The geometric series converges, so also does, by Theorem 6. Thus, we will assume that a = 1. Proof: E(X) = xx Pr(X = x) x E(X)x Pr(X = x) x E(X) E(X) Pr(X = x) = E(X) x E(X) Pr(X = x) = E(X) Pr(X E(X)) Example: If X is B100;1=2, then Pr(X 100) = Pr(X 2E(X)) 1 2 Thisisnotaparticularlyusefulestimate. Note, the disk of convergence ends exactly at the singularity z= 1. Get an answer for 'State and prove Raabe's Test. Note that, is not a geometric series. The difference is that while the Ratio Test for series tells us only that a series converges (ab-solutely), the theorem above tells us that the sequence converges to zero. Does X1 1 2n n3 converge or diverge?. You must use. Test results for water quality (specifically, fecal coliform bacteria concentrations) are sometimes reported as geometric means. Since every converging sequence is bounded, the s n are bounded. Sequences and Series teaches students how to define, notate and interpret different types of series and sequences, such as arithmetic and geometric, and how to use mathematical induction in proofs and on their homework. relate geometric theorems on points, lines, and planes • Logic : Student will use inductive reasoning to draw reasonable conclusions, or deductive reasoning to prove basic theorems, and write conditional statements, converses, inverses and contrapositives. 6Theorem (Alternating series test): (i) Let be an alternating series such that (ii) Then is convergent. The idea is. INDEX proof, 84–86 characteristic polynomial, 292 circle as parametric curve, 259–260 in polar coordinates, 269 closed interval, 11 codomain, 13. Now if 1 k=N+2 a k converged, this implies a k!0. If it converges, then find its sum. Anionic dyes were not used to avoid the potentially confounding. We would like to show you a description here but the site won't allow us. But the integral test easily shows that this series diverges. Proof: Suppose the sequence converges to zero and is monotone decreasing. And hope Now i'm a section of helping you to get a better product. Geometric Sequences - Module 12. then completeness. In these notes we will prove the standard convergence tests and give two tests that aren't in our text. Therefore, the series has bounded partial sums; hence, this sum converges. The reason the test works is that, in the limit, the series looks like a geometric series with ratio L. So, each of the following is geometric. For example, the series. There are several questions Consider the geometric series: (Sum from k=0 to infinity) of ar^k and consider the repeating decimal. A SHORT(ER) PROOF OF THE DIVERGENCE OF THE HARMONIC SERIES LEO GOLDMAKHER It is a classical fact that the harmonic series 1+ 1 2 + 1 3 + 1 4 + diverges. For example, an interesting series which appears in many practical problems in science, engineering, and mathematics is the geometric series + + + + ⋯ where the ⋯ indicates that the series continues indefinitely. THE RATIO TEST We now know how to handle series which we can integrate (the Integral Test), and series which are similar to geometric or p-series (the Comparison Test), but of course there are a. Proof by induction use the basis that if it is true for n = 1 and we assume it is true for n = 1 to some number k and if we can show that it also work for k+1 then we can proof the validity of the geometric progression. Sequences and Series teaches students how to define, notate and interpret different types of series and sequences, such as arithmetic and geometric, and how to use mathematical induction in proofs and on their homework. The nth-term test for divergence. Geometric Proof Video 366. 93Mb) Get help for solving geometric proofs with artificial intelligence methods. If it converges, then find its sum. As an example the geometric series given in the introduction,. Proof: If and are convergent, then it follows from the sum theorem for convergent sequences that is convergent and is valid. This task contains an opportunity to compare the growth of arithmetic and geometric sequences. We begin by giving the following estimate for the partial sum of a p-series: Lemma. Menu Algebra 2 / Sequences and series / Geometric sequences and series A geometric sequence is a sequence of numbers that follows a pattern were the next term is found by multiplying by a constant called the common ratio, r. The series P 1 n=1 1 2 converges (p-series with p= 2 >1). Definition. A geometric series is the sum of terms with a common ratio. For the above proof, using the summation formula to show that the geometric series "expansion" of 0. About This Quiz & Worksheet. But then after the rst N terms the series P nxn is dominated by the geometric series for rjxj, hence converges 8. geometric - definizione, significato, pronuncia audio, sinonimi e più ancora. As a counterexam-ple, few series more clearly illustrate that the convergence of terms. Good GCSE maths result. ALTERNATING SERIES Does an = (−1)nbn or an = (−1)n−1bn, bn ≥ 0? NO Is bn+1 ≤ bn & lim n→∞ YES n = 0? P YES an Converges TELESCOPING SERIES Dosubsequent termscancel out previousterms in the sum? May have to use partial fractions, properties of logarithms, etc. series is a geometric series, our results on geometric series can be used instead. So, each of the following is geometric. Proof: E(X) = xx Pr(X = x) x E(X)x Pr(X = x) x E(X) E(X) Pr(X = x) = E(X) x E(X) Pr(X = x) = E(X) Pr(X E(X)) Example: If X is B100;1=2, then Pr(X 100) = Pr(X 2E(X)) 1 2 Thisisnotaparticularlyusefulestimate. Note: When you talk about an arithmetic sequence, the word arithmetic (in this context) is pronounced air-ith-ME-tic; that is, the accent is on the third syllable. Observe that 1 n2 + 3 < 1 n2 for every n 1. Geometry - MA3110 IC Scope and Sequence Unit Lesson Lesson Objectives Similar Figures Determine if two polygons are similar using dilations. 3 (Part 2) Infinite Geometric Series - Module 12. But if a series converges absolutely, it also converges. Sequences and series, whether they be arithmetic or geometric, have may applications to situations you may not think of as being related to sequences or series. Sequences and Series teaches students how to define, notate and interpret different types of series and sequences, such as arithmetic and geometric, and how to use mathematical induction in proofs and on their homework. It may be one of the most useful tests for convergence. 2 Tests for Convergence Let us determine the convergence or the divergence of a series by comparing it to one whose behavior is already known. The key steps are check cost, condition of pre-order and price recommendation. Its proof is on the separate handout. Geometric Series Convergence. The Limit Superior/Inferior Ratio Test for Series of Complex Numbers. Series with non-negative terms: Comparison test, Cauchy's condensation theorem. Monday, April 30, 2018. Alternating Series test We have the following test for such alternating series: Alternating Series test If the alternating series X1 n=1 ( 1)n 1b n = b 1 b 2 + b 3 b 4 + ::: b n > 0 satis es (i) b n+1 b n for all n (ii) lim n!1 b n = 0 then the series converges. Either the integral test or the Cauchy condensation test shows that the p -series converges for all p > 1 (in which case it is called the over-harmonic series) and diverges for all p ≤ 1. Since , this series diverges. Here, the common ratio (base) is r = sin 2 x , which is always bounded by 1. The sum of two solutions and a scalar multiple of a solution of such a system is again a solution of the system. In Chapter 1, you learned some basic geometric concepts. If it converges, then find its sum. Conditional Convergence; Summary of Tests; Taylor and Maclaurin. of Convergence. 13 GEOMETRIC SERIES, POWER SERIES, RATIO TEST 3 13. N ∞ (−1) j. For example, an interesting series which appears in many practical problems in science, engineering, and mathematics is the geometric series + + + + ⋯ where the ⋯ indicates that the series continues indefinitely. Prove the convergence of the geometric series using $\epsilon$, N definition $\begingroup$ What I gave is the proof you Uniform Convergence using Abel's test. To solve, determine the value of the cumulative distribution function (cdf) for the geometric distribution at x equal to 3. usually shown using the integral test. We've already looked at these. Students may accept the formula for the sum of an infinite geometric series given that $$\left| r \right| < 1$$, and they may even understand the proof of this formula; but they usually are not shown (informally) that the defining feature of a convergent infinite series is that the limit of the series is the limit of its sequence of partial. 19981 THE GEOMETRIC SERIES IN CALCULUS 37--. 1 (A convergent series) A boy is given a chocolate, and he decides that each day he is going to eat. geometric series: Consider the following sequence of numbers: 2, 8, 32, 128, … This is called a geometric progression with a geometric ratio of four. The idea of the ratio test is that we compute the absolute value of that ratio, and. For instance the series 1+3x+9x2+27x3+81x4+ +3kxk+. Proof To prove Property 1, assume that and choose such that By the definition of the limit of a sequence, there exists some such that for all Therefore, you can write the following inequalities. Then is a null sequence, so is a null sequence (by Theorem 7. We must now compute its sum. The Egyptians used this method of finite geometric series mainly to "solve problems dealing with areas of fields and volumes of granaries" but used it for many other uses too, including the pyramids and math problems similar to those one might find on a STAAR test today (see D1, and F1). Print Working with Geometric Sequences Worksheet Recognize a geometric sequence and the quiz allows me to test their knowledge on whatever subject in social studies I am teaching at the. We’ve already looked at these. An infinite geometric series does not converge on a number. You have to to look the same items to test cost as it sometimes will help you in purchasing Design Toscano Set Of 2 Marble Obelisks Geometric Sculptures. Intro to Geometric Sequences. If a n;b n 0 and a n b n for all n, then if the series X1 n=1 b n converges then the series X1 n=1 a n converges. A note about the geometric series Before we get into today's primary topic, I have to clear up a little detail about the geometric series. It's pretty simple; we're really just asking whether the n th term of a series converges to zero, but the divergence test has some important limitations that we should get out of the way right away. Infinite series and the biggest maths problem of them all: One famous series is the Riemann zeta function, which is involved in one of the biggest open problems in maths: the Riemann hypothesis. The term r is the common ratio, and a is the first term of the series. Most importantly, it also tells you the areas which need your attention. Test for Divergence: If lim → ≠0, then diverges. Where a 1 = the first term, a 2 = the second term, and so on a n = the last term (or the n th term) and a m = any term before the last term. Rather than write down the proof, let me just motivate why this theorem is true. for some constants a and r. Then, once you get an explicit formula for f ( x ), you can plug in x = π/3. Another way of looking at that is to ask for an average number of trials before the first occurrence of the event. The Harmonic Series Diverges Again and Again∗ Steven J. Finally, we give a test which helps us to analyze convergence of an alternating series. P), then b is the arithmetic mean of a and c. The Ratio Test Proof (1): If 0 ˆ<1, we can apply the previous theorem to see P 1 n=1 ja nj converges. A proof is a mathematical argument used to verify the truth of a statement. The geometric series diverges if. (Notethatin thissectionwewill sometimesbeginourseriesat n 0 and sometimesbegin them at n 1. For this question you will write a two-column proof of the first part of the Overlapping Angle Theorem. Proof of 1 (if L < 1, then the series converges) Proof of 2 (if L > 1, then the series diverges) Proof of 1 (if L < 1, then the series converges) Our aim here is to compare the given series with a convergent geometric series (we will be using a comparison test). The common ratio of the series is positive. The key steps are check cost, condition of pre-order and price recommendation. Determine whether the series converges or diverges. ALTERNATING SERIES Does an = (−1)nbn or an = (−1)n−1bn, bn ≥ 0? NO Is bn+1 ≤ bn & lim n→∞ YES n = 0? P YES an Converges TELESCOPING SERIES Dosubsequent termscancel out previousterms in the sum? May have to use partial fractions, properties of logarithms, etc. ap calculus bc college board calculus nth term and geometric series test fun worksheet. N ∞ (−1) j. Complete the Lesson Practice homework help geometry proofs homework. $\begingroup$ I agree that the link to geometric series is the only compelling reason to consider writing the finite geometric series formula in the "weird" way. On The Ratio Test for Positive Series of Real Numbers we looked at a very useful test for determining the convergence of a series of real numbers. Peterson Department of Biological Sciences and Department of Mathematical Sciences Clemson University February 1, 2019 Outline Geometric Series The Ratio Test The Root Test Examples A Taste of Power Series. Theorem (A Divergence test): If the series is convergent, then The test for divergence: If denotes the sequence of partial sums of then if does not exist or if , then the series is divergent. The value to which the series converges is the least of all possible upper bounds. If L= 1 then the test fails. Proof: The convergence properties of the power series are a consequence of the ratio test. Test results for water quality (specifically, fecal coliform bacteria concentrations) are sometimes reported as geometric means. The Mathematics Vision Project: Scott Hendrickson, Joleigh Honey, Barbara Kuehl, Travis Lemon, Janet Sutorius. (b) (i) Matrices over R, The matrix representation of systems of homogeneous and non- homogeneous linear equations. 4 Ratio test The geometric series leads to a useful test for convergence of the general series X1 n=0 a n= a 0 + a 1 + a 2 + (12) We can make sense of this series again as the limit of the partial sums S n = a 0 + a 1 + + a n as n!1. There are several questions Consider the geometric series: (Sum from k=0 to infinity) of ar^k and consider the repeating decimal. P), then b is the arithmetic mean of a and c. The common ratio in one geometric sequence is a whole number and in the other sequence it is a percent. Geometric Distribution De nition (Geometric Distribution) In a series of Bernoulli trials (independent trials with constant probability p of success), let the random variable Xdenote the number of trials until the rst success. The geometric series converges, and so, by the Direct Comparison Test, the. Geometric Series and the Test for Divergence - Part 1 patrickJMT. com where we believe that there is nothing wrong with being square! This page includes Geometry Worksheets on angles, coordinate geometry, triangles, quadrilaterals, transformations and three-dimensional geometry worksheets. Water Quality Standards. How is this done? Sn=a1(1-r^n)/(1-r). Note: If lim → =0, then there is no conclusion about. Peterson Department of Biological Sciences and Department of Mathematical Sciences Clemson University February 1, 2019 Outline Geometric Series The Ratio Test The Root Test Examples A Taste of Power Series. Sequences and series, whether they be arithmetic or geometric, have may applications to situations you may not think of as being related to sequences or series. 3 (Part 3) Review of Module 12 on Sequences and Series. 21) a 1 = −2, r = 5, S n = −62 22) a 1 = 3, r = −3, S n = −60 23) a 1 = −3, r = 4, S n = −4095 24) a 1 = −3, r = −2, S n = 63 25) −4 + 16 − 64 + 256 , S n = 52428 26) Σ m = 1 n −2 ⋅ 4m − 1 = −42-2-. Real and Integer Numbers Calculators and Percentages Sum of Positive Integers Calculator. Get the free "Infinite Series Analyzer" widget for your website, blog, Wordpress, Blogger, or iGoogle. The shape of the material used for the lampshade is a sector of a circle. Geometric Sequences - Module 12. What I want to do in this video is now think about the sum of an infinite geometric series. lcprefrigeration. If 1 then 1 for all so does not define a null sequence and the series diverges by the null sequence test. The circle has radius 25 centimetres and the angle of the sector is 280° a) Find the area of the sector of the circle. 4 Ratio test The geometric series leads to a useful test for convergence of the general series X1 n=0 a n= a 0 + a 1 + a 2 + (12) We can make sense of this series again as the limit of the partial sums S n = a 0 + a 1 + + a n as n!1. Proof - Convergence of a Geometric Series Contact Us If you are in need of technical support, have a question about advertising opportunities, or have a general question, please contact us by phone or submit a message through the form below. Def: The geometric series is convergent if |r. These online tests are designed to work on computers, laptops, iPads, and other tablets. The geometric mean isn’t affected by those factors. Then the series is called Geometric Series. We will now look at a more general root test which can be applied to series of complex numbers. Peterson Department of Biological Sciences and Department of Mathematical Sciences Clemson University February 1, 2019 Outline Geometric Series The Ratio Test The Root Test Examples A Taste of Power Series. An arithmetic sequence is a sequence with the difference between two consecutive terms constant. In fact, our proof is an extension of the nice result given by Cohen and Knight [2]. (b) If , the series diverges. To solve problems on this page, you should be familiar with arithmetic progressions geometric progressions arithmetic-geometric progressions. The series , which was one of our examples given above, is a geometric series since =. - In a paragraph proof, statements and their justifications are written in sentences in a logical order. For example, the series. C2 Sequences & Series - Arithmetic & Geometric Series 3 MS C2 Sequences & Series - Arithmetic & Geometric Series 3 QP C2 Sequences & Series - Arithmetic & Geometric Series 4 MS. Since , this series diverges. The Limit Superior Root Test for Series of Complex Numbers. In a Geometric Sequence each term is found by multiplying the previous term by a constant. I offer this post for the enrichment of talented Precalculus students who have exhibited mastery of geometric series. 4 Ratio test The geometric series leads to a useful test for convergence of the general series X1 n=0 a n= a 0 + a 1 + a 2 + (12) We can make sense of this series again as the limit of the partial sums S n = a 0 + a 1 + + a n as n!1. You can see that this is reasonable by dividing 1 by , or using the the formula for the sum of a geometric series with ratio. Test for Divergence: If lim → ≠0, then diverges. The key steps are check cost, condition of pre-order and price recommendation. mathematicsvisionproject. ' and find homework help for other Math questions at eNotes If L=1 the test is inconclusive. org right now: https://www. The sum of a convergent geometric series can be calculated with the formula a ⁄ 1-r , where “a” is the first term in the series and “r” is the number getting raised to a power. Mathematical Induction - Problems With Solutions Several problems with detailed solutions on mathematical induction are presented. Geometric Series The formula for the sum of the first n terms of a geometric series is derived by using several ideas, each expressed concisely with subscripts and exponents. As this was clearly the case for the geometric series ∑ n = 0 ∞ x n (for 0 ≤ x 1), he asserted (perhaps as a reminder of Euclid's proof for the area of a circle) that a series ∑ n = 0 ∞ a n with positive terms will converge if, for n large enough, a n +1 is at most 1 2 a n; and without bothering about details, he stated correctly. Since the proof of the ratio test relies on this convergence, it is circular to argue that a geometric series converges by the ratio test (unless, of course, you have another proof for the ratio test that doesn't use the convergence of geometric series). In fact this series diverges quite slowly. A telescoping series is any series where nearly every term cancels with a preceeding or following term. ALTERNATING SERIES Does an = (−1)nbn or an = (−1)n−1bn, bn ≥ 0? NO Is bn+1 ≤ bn & lim n→∞ YES n = 0? P YES an Converges TELESCOPING SERIES Dosubsequent termscancel out previousterms in the sum? May have to use partial fractions, properties of logarithms, etc. Check out the activities offered on this site. Show that if X1 k=0 a k converges and b k is a bounded sequence, then X1 k=0 a kb k. Proof: If and are convergent, then it follows from the sum theorem for convergent sequences that is convergent and is valid. $\begingroup$ I agree that the link to geometric series is the only compelling reason to consider writing the finite geometric series formula in the "weird" way. We would like to show you a description here but the site won't allow us. Arithmetic and geometric series are two the most simple series for calculating the sum of numbers in series in a very simple way. Introduction to Video: Proof by Mathematical Induction. The comparison test allows us to construct other examples from this: THEOREM 2 (Comparison Test). A geometric series is a). To find the sum of an infinite geometric series having ratios with an absolute value less than one, use the formula, S = a 1 1 − r , where a 1 is the first term and r is the common ratio. 1 Sequences We call a list of numbers written down in succession a sequence; for example, the numbers drawn in a lottery: 12,22,5,6,16,43, For this sequence, there is no clear rule that will enable you to predict with certainty the next number in the sequence. We begin by giving the following estimate for the partial sum of a p-series: Lemma. For general help, questions, and suggestions, try our dedicated support forums. A geometric series is a). By inspection, it can be difficult to see whether a series will converge or not. Note: If lim → =0, then there is no conclusion about. the "Ratio Test for a Sequence". then completeness. Theorem 2: If the power series f(x) = P∞ n=0 anx n is convergent at x = R, then it is a continuous function within the interval of convergence. Suppose the interest rate is loo%, so i = 1. Sum of Arithmetic Geometric Sequence In mathematics, an arithmetico-geometric sequence is the result of the term-by-term multiplication of a geometric progression with the corresponding terms of an arithmetic progression. Writing a proof can even be more daunting I kept the reader(s) in mind when I wrote the proofs outlines below. I have to calculate the standard 90% confidence intervals of the ratio test/reference (T/R) and I have to answer to this question: the products were considered bioequivalent if the difference between two. That is, if we can prove that the sequence {a n} does not. This ratio is called the common ratio. The following rules are often helpful:. In mathematical analysis, the alternating series test is the method used to prove that an alternating series with terms that decrease in absolute value is a convergent series. Here it is. As a counterexam-ple, few series more clearly illustrate that the convergence of terms. The Ratio Test Proof (1): If 0 ˆ<1, we can apply the previous theorem to see P 1 n=1 ja nj converges. The sum of the first n terms of the geometric sequence, in expanded form, is as follows:. can inductive. Then there is a such that for we have. The geometric sequence can be rewritten as where is the amount of terms, is the common ratio, and is the first term. 1 The Geometric Series (page 373) EXAMPLE. How do we distinguish graphically between an arithmetic and a geometric sequence? 9. Infinite series. We would like to show you a description here but the site won't allow us. Introduction to Video: Proof by Mathematical Induction. Since lim n→∞ n √ a n = ρ < 1, then for any > 0, small enough such that ρ + = r < 1, there exists N large with. ALTERNATING SERIES Does an = (−1)nbn or an = (−1)n−1bn, bn ≥ 0? NO Is bn+1 ≤ bn & lim n→∞ YES n = 0? P YES an Converges TELESCOPING SERIES Dosubsequent termscancel out previousterms in the sum? May have to use partial fractions, properties of logarithms, etc. Say we have an infinite geometric series whose first term is aaaa and common ratio is rrrr. A SHORT(ER) PROOF OF THE DIVERGENCE OF THE HARMONIC SERIES LEO GOLDMAKHER It is a classical fact that the harmonic series 1+ 1 2 + 1 3 + 1 4 + diverges. 93Mb) Get help for solving geometric proofs with artificial intelligence methods. Any one of these nite partial sums exists but the in nite sum does not necessarily converge. Lecture 25/26 : Integral Test for p-series and The Comparison test In this section, we show how to use the integral test to decide whether a series of the form X1 n=a 1 np (where a 1) converges or diverges by comparing it to an improper integral. Hey guys, I'm stuck on wording of a homework assignment and thought you might be able to help me. A geometric sequence is a sequence with the ratio between two consecutive terms constant. The sequence 16 ,8 ,4 ,2 ,1 ,1/2 ,… = is a decreasing geometric sequence of common ratio ½. Infinite series can be daunting, as they are quite hard to visualize. can inductive. For such series, it necessary to evaluate limits of th roots of more complicated exp8 ressions. 832 respectively. Since every converging sequence is bounded, the s n are bounded. Grappling with the geometric series, geometry formulas or geometric sequence? Our tutors can help. Compute Geometric Distribution cdf. The term r is the common ratio, and a is the first term of the series. This usually takes the form of a formal proof, which is an orderly series of statements based upon axioms, theorems, and statements derived using rules of inference. Alternating Series test We have the following test for such alternating series: Alternating Series test If the alternating series X1 n=1 ( 1)n 1b n = b 1 b 2 + b 3 b 4 + ::: b n > 0 satis es (i) b n+1 b n for all n (ii) lim n!1 b n = 0 then the series converges. If it converges, then find its sum. For those geometric series, the signs of the terms alternate between positive and negative. In mathematical analysis, the alternating series test is the method used to prove that an alternating series with terms that decrease in absolute value is a convergent series. org/math/precalculus/seq_induction/infinite-geometric-series/e/understand. 9is formulated for convergent series, its main importance is as a \divergence test": if the general term in an in nite series does not tend to 0 then the series diverges. understand the concept of a geometric series • use and manipulate the appropriate formula • apply their knowledge of geometric series to everyday applications • deal with combinations of geometric sequences and series and derive information from them • find the sum to infinity of a geometric series, where -1 < r < 1. Interpret the structure of expressions. Convergence & Divergence - Geometric Series, Telescoping Series, Harmonic Series, Divergence Test - Duration: 50:43. Finally, we give a test which helps us to analyze convergence of an alternating series. We will now look at a more general root test which can be applied to series of complex numbers. Properties of Series; Arithmetic Series; Finite Geometric Series; Infinite Geometric Series; Decimal Expansion; Word Problems; Visualization of Series; The Divergence Test; The Alternating Series Test; The Ratio Test; The Integral Test; The Comparison Test; Absolute Convergence vs. The standard proof involves grouping larger and larger numbers of consecutive terms,. Where a 1 = the first term, a 2 = the second term, and so on a n = the last term (or the n th term) and a m = any term before the last term. This looks to me like an expression of the "ratio test" for convergence of a series. series is a geometric series, our results on geometric series can be used instead. In the following series, the numerators are in AP and the denominators are in GP:. Stamps Prairie State College The harmonic series, X∞ n=1 1 n = 1+ 1 2 + 1 3 + 1 4 + 1 5 +···, is one of the most celebrated infinite series of mathematics. If the limit of a[n] is not zero, or does not exist, then the sum diverges. Finally, we give a test which helps us to analyze convergence of an alternating series. Determine whether the series converges or diverges. Infinite series can be daunting, as they are quite hard to visualize. Absolute Convergence If the series |a n | converges, then the series a n also converges. Our approach can be described as follows. However, notice that both parts of the series term are numbers raised to a power. 717171717171 for these problems: Question 1: Find a formula for the n-th partial sum of the series and PROVE your result using the Cauchy Convergence Criterion. Select one of the links below to get started. Does X1 1 2n n3 converge or diverge?. Infinite Sum Geometric Series. Let a n a n+1 0. | 2020-04-07T23:41:06 | {
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https://math.stackexchange.com/questions/2188962/determine-if-this-series-converges-or-diverges | # Determine if this series converges or diverges
$$\sum_{n=1}^\infty \sin^{[n]}(1)$$ Where by $\sin^{[n]}(1)$ we mean $\sin\left(\sin\left(\dots\sin(1)\right)\right)$ composed $n$ times.
Have tried the divergence test, which fails. Have tried Ratio test, also fails, as the limit is 1. Integral test, or root test do not seem promising. Help is appreciated
• There are higher order ratio tests that may be useful (I have no clue) – Mark Mar 16 '17 at 5:52
• math.stackexchange.com/questions/14004/… – Michael Biro Mar 16 '17 at 5:54
• @Mark Ratio test will return one, hence, it is inconclusive. :P – Simply Beautiful Art Mar 16 '17 at 13:07
• Actually, the initial sine can be of anything real because it will always be in [-1,1]. This might be more interesting if we allow a complex sine. – richard1941 Mar 22 '17 at 22:00
• @Mark: I was unable to apply Raabe's Test successfully. – robjohn Feb 3 at 16:57
Once we prove the inequality $$\sin x > \frac{x}{1+x} \qquad (\star)$$ for $x \in (0,1]$, we can inductively show that $\sin^{[n]}(1) > \frac1n$ when $n\ge 2$. We have $\sin \sin 1 \approx 0.75 > \frac12$, and whenever $\sin^{[n]}(1) > \frac1n$, we have $$\sin^{[n+1]}(1) = \sin \sin^{[n]}(1) > \sin \frac1n > \frac{1/n}{1+1/n} = \frac1{n+1}.$$ Therefore, by the comparison test, $$\sum_{n=1}^\infty \sin^{[n]}(1) = \sin 1 + \sum_{n=2}^\infty \sin^{[n]}(1) > \sin 1 + \sum_{n=2}^\infty \frac1n,$$ which diverges.
To prove $(\star)$... well, to be honest, I just graphed both sides. But we can prove that $\sin x > x - \frac{x^3}{6}$ on the relevant interval by thinking about the error term in the Taylor series, and $x - \frac{x^3}{6} > \frac{x}{1+x}$ can be rearranged to $(x+1)(x -\frac{x^3}{6}) - x > 0$, which factors as $-\frac16 x^2(x-2)(x+3) > 0$.
• Let $f(x) = (x+1)\sin(x)-x$. We want to show $f(x) \ge 0$. Note $f(0) = 0$ and $f'(x) = \sin(x)+(x+1)\cos(x)-1$. It suffices to show $f'(x) \ge 0$. So it suffices to show $\sin(x)+\cos(x) \ge 1$ on $[0,1]$ (since $x\cos(x) \ge 0$). But by multiplying by $\frac{\sqrt{2}}{2}$, this is equivalent to $\sin(x+\frac{\pi}{4}) \ge \frac{\sqrt{2}}{2}$. And this is true for $x \in [0,\frac{\pi}{2}]$, so we're done since $\frac{\pi}{2} > 1$. – mathworker21 Mar 16 '17 at 6:07
Another possibility is to try to find an equivalent of the general term of this sequence : $\begin{cases} a_0=1\\ a_{n+1}=sin(a_n) \end{cases}$
Note that $f(x)=sin(x)$ has derivative $f'(x)=cos(x)$ which is positive on $[0,a_0]$ and also $<1$ on $]0,a_0[$ so $f$ is a contraction. From there it is easy to prove that $a_n\to 0$.
This means that $a_{n+1}\sim a_n$ when $n\to\infty$.
In this kind of problem we always search for an $\alpha$ such that $\mathbf{(a_{n+1})^\alpha-(a_n)^\alpha}$ does not depend of $\mathbf{a_n}$ (in the $\sim$ sense of the expression) so we are able to solve the recurrence.
From Taylor expansion $\displaystyle{(a_{n+1})^\alpha = \bigg(a_n - \frac{a_n^3}{6}+o(a_n^4)\bigg)^\alpha}=(a_n)^\alpha\bigg(1 - \frac{a_n^2}{6}+o(a_n^3)\bigg)^\alpha=(a_n)^\alpha\big(1-\frac{\alpha}{6}a_n^2+o(a_n^3)\big)$
So $(a_{n+1})^\alpha-(a_n)^\alpha=-\frac{\alpha}{6}(a_n)^{\alpha+2}+o((a_n)^{\alpha+3})\quad$ we see that we need $\alpha=-2$
Let's put $b_n=1/(a_n)^2,\qquad$ $b_n\to\infty$
We have $b_{n+1}-b_n=\frac 13+o(1/b_n)$ thus $b_n\sim\frac n3\qquad$
(more precisely $b_n=n/3+o(\ln(n))$ but it is not important at this point).
Finally $a_n\sim\sqrt\frac 3n,$ which is a term of a divergent serie so $\sum a_n$ diverges as well.
• This is also very nice! Thank you ! – userX Mar 19 '17 at 2:57
I wrote this answer for another question, but this question was pointed out in a comment, so I posted it here.
$$|\,a_{n+1}|=|\sin(a_n)\,|\le|\,a_n|$$. Thus, $$|\,a_n|$$ is decreasing and bounded below by $$0$$. Thus, $$|\,a_n|$$ converges to some $$a_\infty$$, and we must then have $$\sin(a_\infty)=a_\infty$$, which means that $$a_\infty=0$$.
Using $$\sin(x)=x-\frac16x^3+O\!\left(x^5\right)$$, we get \begin{align} \frac1{a_{n+1}^2}-\frac1{a_n^2} &=\frac1{\sin^2(a_n)}-\frac1{a_n^2}\\ &=\frac{a_n^2-\sin^2(a_n)}{a_n^2\sin^2(a_n)}\\ &=\frac{\frac13a_n^4+O\!\left(a_n^6\right)}{a_n^4+O\!\left(a_n^6\right)}\\ &=\frac13+O\!\left(a_n^2\right)\tag1 \end{align} Stolz-Cesàro says that $$\lim_{n\to\infty}\frac{\frac1{a_n^2}}n=\frac13\tag2$$ That is, $$\bbox[5px,border:2px solid #C0A000]{a_n\sim\sqrt{\frac3n}}\tag3$$ which means that the series diverges.
Motivation for $$\boldsymbol{(1)}$$
Note that $$a_{n+1}-a_n=\sin(a_n)-a_n\sim-\frac16a_n^3\tag4$$ which is a discrete version of $$\frac{\mathrm{d}a_n}{a_n^3}=-\frac{\mathrm{d}n}6\tag5$$ whose solution is $$\frac1{a_n^2}=\frac{n-n_0}3\tag6$$ so that $$\frac1{a_{n+1}^2}-\frac1{a_n^2}=\frac13\tag7$$ which suggests $$(1)$$. | 2020-05-26T02:58:05 | {
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https://math.stackexchange.com/questions/2545658/counting-with-recurrence-strings-which-avoid-the-substring-ab | # Counting with recurrence: Strings which avoid the substring $AB$.
Let $X(n)$ be the number of length $n$ strings from the alphabet $\{A,B,C,D\}$ which avoid the substring $AB$. Write a recurrence for $X(n)$ and explain in words why it satisfies the recurrence.
My attempt:
Any length $n$ string sequence must start with either $A,B,C,D$. If it starts with $B,C,D$, we don't need to worry about the what follows right after, so the number of length $n$ strings that avoid $AB$ and start with $B,C,D$ is $3X(n-1)$. Now we deal with the case when it starts with $A$. If it starts with $A$, the next letter must be either $A,C,D$, only three choices, which means the rest of the sequence can be any $n-2$ length string which avoids $AB$. The number of length $n$ strings which start with $A$ and which avoid $AB$ is $3X(n-2)$, so we get \begin{cases} X(0)=1\\[4pt] X(1)=4\\[4pt] X(n)=3X(n-1)+3X(n-2),\;\text{if}\;n\ge 2\\ \end{cases} However the solution is: \begin{cases} X(0)=1\\[4pt] X(1)=4\\[4pt] X(n)=4X(n-1)-X(n-2),\;\text{if}\;n\ge 2\\ \end{cases} Note that for $n=2$, the results are the same, but they are not the same for $n=3$.
Please let me know what I missed in my attempt. Thank you.
As to where you went wrong, your statement
$\;\;\;\;$The number of length $n$ strings which start with $A$ and which avoid $AB$ is $3X(n-2)$
is not correct, since for $n \ge 3$, it includes some strings which start with $AAB$.
A correct recursion can be achieved as follows . . .
By direct count, we get $x(0) = 1$ (the empty string), and $x(1) = 4$ (any string of length $1$).
Suppose $n \ge 2$.
Consider the count of all $n$-term sequences $s$ for which the $(n-1)$-term subsequence $s'$ with the first term of $s$ removed is such that $s'$ avoids the subsequence $AB$.
Since we are allowing all possible first terms for $s$, the count is $4x(n-1)$.
But this is an overcount since it allows sequences $s$ which start with $AB$.
To correct the count, subtract the count of those already counted sequences which start with $AB$.
Thus, we need to subtract the count of the sequences $s=ABs''$ such that the $(n-2)$-term sequence $s''$ avoids the subsequence $AB$, and there are $x(n-2)$ of those.
Hence, for $n \ge 2$, we get $x(n) = 4x(n-1) - x(n-2)$.
As an alternate approach, you can use a linked recursion.
It's a little more work, but I think the logic is simpler.
Here's how it would go . . .
Let $x(n)$ be the number of qualifying $n$-term sequences with no immediately preceding $A$, and let $y(n)$ be the number of qualifying $n$-term sequences with an immediately preceding $A$.
Then for $n > 0$, we get \begin{align*} x(n) &= 3x(n-1) + y(n-1)\\[4pt] y(n) &= y(n-1)+ 2x(n-1) \end{align*} and for $n=0$, we have $x(0) = 1,\;y(0) = 1$.
To unlink the recursion, solve the equation $$x(n) = 3x(n-1) + y(n-1)$$ for $y(n-1)$, which yields $$y(n-1) = x(n) - 3x(n-1)$$ which implies $$y(n) = x(n+1) - 3x(n)$$ Then, replacing $y(n)$ and $y(n-1)$ in the equation $$y(n) = y(n-1)+ 2x(n-1)$$ yields $$x(n+1) = 4x(n) - x(n-1)$$ or equivalently, $$x(n) = 4x(n-1) - x(n-2)$$ valid for all $n > 1$, with initial values $x(0) = 1,\;x(1) = 4$.
In his answer quasi has pointed out your error. As for the correct recurrence, you can argue as follows:
An admissible word of length $n$ begins with one of $A$, $B$, $C$, $D$ and is followed by an admissible word of length $n-1$. We therefore would have $X(n)=4 X(n-1)$. Now this does include words which begin with $AB$, but are otherwise admissible. There are $X(n-2)$ of these. Discounting them we obtain the recursion $$X(n)=4X(n-1)-X(n-2)\ .$$ | 2019-11-19T12:34:20 | {
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http://mathhelpforum.com/calculus/170077-trig-limit-when-numerator-undefined.html | # Math Help - Trig Limit when the Numerator is Undefined
1. ## Trig Limit when the Numerator is Undefined
Can somebody help me understand this one, please?
$\displaystyle \lim_{x\to \infty } \frac{\cos^4(x)}{5 + x^2}$
I get that the answer is zero, because the denominator goes to infinity due to the $x^2$. I also understand that the value of $\cos^4(x)$ is always between 0 and 1, which is why we can conclude that the limit overall is zero.
What is a little weird for me is that the limit of $\cos^4(x)$ doesn't exist. I can't get past the idea that the limit of a fraction can be defined when the limit of the numerator doesn't exist. What if the numerator was simply $\cos (x)$? Then what would the overall limit be (since $\cos (x)$ can vary between -1 and 1)?
Sorry if I'm being a little dense, but can somebody try to explain this to me?
Thanks.
2. Originally Posted by joatmon
Can somebody help me understand this one, please?
$\displaystyle \lim_{h\to \infty } \frac{\cos^4(x)}{5 + x^2}$
I get that the answer is zero, because the numerator goes to infinity due to the $x^2$. I also understand that the value of $\cos^4(x)$ is always between 0 and 1, which is why we can conclude that the limit overall is zero.
What is a little weird for me is that the limit of $\cos^4(x)$ doesn't exist. I can't get past the idea that the limit of a fraction can be defined when the limit of the numerator doesn't exist. What if the numerator was simply $\cos (x)$? Then what would the overall limit be (since $\cos (x)$ can vary between -1 and 1)?
Sorry if I'm being a little dense, but can somebody try to explain this to me?
Thanks.
The idea is hidden in what you have said. The squeeze theorem gives you what you need. You have said that you believe that
$\displaystyle -1 \le \cos(x) \le 1$ for all x. Now if we divide this inequality by $x^2+1$ Since this is always non negative we don't have to change the inequality signs we get
$\displaystyle \frac{-1}{x^2+1} \le \frac{\cos(x)}{x^2+1} \le \frac{1}{x^2+1}$
This inequality is still valid for all values of x
Now if we take the limit with out knowing what the middle limit its we will know two numbers that it is still in between by the inequalities.
$\displaystyle \lim_{x \to \infty}\frac{-1}{x^2+1} \le \lim_{x \to \infty}\frac{\cos(x)}{x^2+1} \le \lim_{x \to \infty}\frac{1}{x^2+1}$
This gives
$\displaystyle 0 \le \lim_{x \to \infty}\frac{\cos(x)}{x^2+1} \le 0$
This tells us that the limit is bigger than $0$, but also less than $0$ so the limit must be zero.
3. I assume this should actually be $\displaystyle \lim_{x \to \infty}\frac{\cos^4{x}}{5 + x^2}$
I think you can squeeze it...
$\displaystyle \frac{0}{5 + x^2} \leq \frac{\cos^4{x}}{5 + x^2} \leq \frac{1}{5 + x^2}$, so
$\displaystyle \lim_{x \to \infty} 0 \leq \lim_{x \to \infty}\frac{\cos^4{x}}{5 + x^2} \leq \lim_{x \to \infty}\frac{1}{5 + x^2}$
$\displaystyle 0 \leq \lim_{x \to \infty}\frac{\cos^4{x}}{5 + x^2} \leq 0$.
So $\displaystyle \lim_{x \to \infty} \frac{\cos^4{x}}{5 + x^2} = 0$.
4. Thanks. I had played with trying to squeeze this out, but couldn't get it to work.
So, my logic about why the limit is zero isn't correct. I had thought that, since the function $cos^4(x)$ is always between 0 and 1, which we know even without using the squeeze theorem, what forces the limit to zero is that the denominator becomes infinitely large.
But what I think that you are saying is that the denominator has nothing to do with it. It's that $\lim_{x\to \infty } \cos^4(x)$ can be shown to be zero by using the squeeze theorem. Thus, no matter what happens to the denominator (so long as its limit exists and is not zero), the value of the limit is still zero since the limit of the numerator is zero.
Is a correct interpretation?
5. Just saw Prove It's reply. Now I get it entirely. Thanks to both of you.
6. Originally Posted by TheEmptySet
The idea is hidden in what you have said. The squeeze theorem gives you what you need. You have said that you believe that
$\displaystyle -1 \le \cos(x) \le 1$ for all x. Now if we divide this inequality by $x^2+1$ Since this is always non negative we don't have to change the inequality signs we get
$\displaystyle \frac{-1}{x^2+1} \le \frac{\cos(x)}{x^2+1} \le \frac{1}{x^2+1}$
This inequality is still valid for all values of x
Now if we take the limit with out knowing what the middle limit its we will know two numbers that it is still in between by the inequalities.
$\displaystyle \lim_{x \to \infty}\frac{-1}{x^2+1} \le \lim_{x \to \infty}\frac{\cos(x)}{x^2+1} \le \lim_{x \to \infty}\frac{1}{x^2+1}$
This gives
$\displaystyle 0 \le \lim_{x \to \infty}\frac{\cos(x)}{x^2+1} \le 0$
This tells us that the limit is bigger than $0$, but also less than $0$ so the limit must be zero.
Having the fourth power of cosine, I'd use 0 as the lower bounds, since the numerator and denomiator are positive.
I see that Profit said the same thing.
7. Originally Posted by joatmon
Thanks. I had played with trying to squeeze this out, but couldn't get it to work.
So, my logic about why the limit is zero isn't correct. I had thought that, since the function $cos^4(x)$ is always between 0 and 1, which we know even without using the squeeze theorem, what forces the limit to zero is that the denominator becomes infinitely large.
But what I think that you are saying is that the denominator has nothing to do with it. It's that $\lim_{x\to \infty } \cos^4(x)$ can be shown to be zero by using the squeeze theorem.
No, it doesn't. The constant function f(x)= 1/2 is always between 0 and 1- does that "force" its limit to be 0?
$\lim_{x\to\infty} cos^4(x)$ does not exist.
Thus, no matter what happens to the denominator (so long as its limit exists and is not zero), the value of the limit is still zero since the limit of the numerator is zero.
Is a correct interpretation?
No, it isn't. It is the fact that the denomator goes to infinity that makes the limit 0. The only thing required of the numerator is that it NOT go to infinity or negative infinity. | 2016-05-25T19:58:32 | {
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https://web2.0calc.com/questions/help_2802 | +0
# help
0
59
13
If $a, b, c$ are positive integers less than 13 such that
$2ab+bc+ca \equiv 0\pmod{13}$
$ab+2bc+ca \equiv 6abc\pmod{13}$
$ab+bc+2ca \equiv 8abc\pmod {13}$
then determine the remainder when $a + b + c$ is divided by $13$.
Sep 27, 2020
#1
0
a + b + c leaves a remainder of 8 when divided by 13.
Sep 27, 2020
#2
0
Sorry, but it says that is incorrect.
Thank you for trying to help out though!
Also, if someone wants to post a solution to this, it'd be best if they provided an explanation, because I'd like to learn from it. Thanks!
Guest Sep 27, 2020
#3
+1
Look at many, many solutions here:
https://www.wolframalpha.com/input/?i=%282ab%2Bbc%2Bca%29+mod+13+%3D0%2C++%28ab%2B2bc%2Bca%29+mod+13%3D+6abc%2C++%28ab%2Bbc%2B2ca%29++mod+13+%3D+8abc%2C+solve+for+a%2Cb%2Cc
Sep 27, 2020
#4
0
Unfortunately, this also does not work as your input had equal signs instead of modular congruence signs. I don't believe you can input modular congruence signs into Wolfram Alpha.
Thanks for trying to help me though!
Guest Sep 28, 2020
#5
+111124
0
If a, b, and c are all 0 then these three equations are all true.
And if this is true then (a+b+c) mod13 = 0
So if there is only one solution then that solution must be 0
However, just two of them are 0 then the third one can equal to many different numbers.
So
I do not believe this has one specific answer.
Sep 28, 2020
#6
+1
I know I must be annoying for saying all of the answers are wrong, but unfortunately there's a problem with this one too. The problem states that $a, b, c$ must be positive integers less than 13, so $0$ doesn't fit this requirement. However, thank you very much for trying to solve this for me!
Guest Sep 28, 2020
#7
+111124
0
Yep, good point!
Melody Sep 29, 2020
#8
0
Sep 29, 2020
#9
+1
a = 3, b = 6, c = 9.
Unique.
Will post a solution, of sorts, if anyone is interested.
Guest Sep 30, 2020
#10
+111124
0
Thanks guest,
I have some interest.
Did you use code to do the grunt work and just test all the possibilities
Or did you actually come by the answer in a logical, mathematical (no trial and error) way ?
Melody Sep 30, 2020
#11
+1
Thanks Melody, it's nice to know that someone is actually interested.
Often that doesn't appear to be the case.
The equations can be written in the form
$$\displaystyle 2X+Y+Z\equiv 0 \dots\dots(1)\\ X+2Y+Z\equiv R\dots\dots(2)\\ X+Y+2Z\equiv S\dots\dots(3)$$
where each one is mod 13.
$$\displaystyle (1) - 2(3) : -Y-3Z \equiv -2S\dots\dots(4)\\ (2) - (3) :Y - Z \equiv R-S\dots\dots(5)\\ (4)+(5) : -4Z\equiv R-3S\dots\dots(6).$$
Returning to the original unknowns,
$$\displaystyle -4ca \equiv 6abc-24abc \equiv -18abc,$$
from which
$$\displaystyle9b=2+13k\quad \text{ where k is an integer}.$$
The smallest value of k for which b will be an integer is k = 4, leading to 9b = 54, b = 6.
The next suitable value for k puts b out of the permitted range.
The remaining unknowns (a = 3 and c = 9) can then be found by back substitution, (6) into (5) and then into (2) or (3).
Leave that to you.
Guest Sep 30, 2020
#12
0
THANK YOU THANK YOU THANK YOU!
I really like your explanation! Thanks a lot!
Guest Sep 30, 2020
#13
+111124
0
Brilliant, thanks very much!
I never would have thought to do those substitutions!
I continued with matrices, and again, I would not have thought to use matrices on a modulo question if you had not inspired me.
Melody Sep 30, 2020 | 2020-10-22T12:57:51 | {
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https://math.stackexchange.com/questions/2612170/computing-convolution-using-the-fourier-transform | # Computing convolution using the Fourier transform
I am given that the functions $x(t)$ and $h(t)$ are defined by $$x(t) = t\, e^{-2t}\, u(t) \qquad \text{and} \qquad h(t) = e^{-4t}\, u(t)$$ where $u(t)$ denotes the unit step function $$u(t) = \begin{cases} 1 & \text{if t \geq 0}\\ 0 & \text{otherwise}. \end{cases}$$
The question is the following:
Compute the convolution $(x * h)(t)$ by finding the corresponding Fourier transform $X(\omega)$ and $H(\omega)$ using the convolution property, and then inverse transforming.
By the convolution property, we mean $(x*h)(t) \stackrel{\text{FT}}{\longleftrightarrow} X(\omega) H(\omega)$.
Now I found the Fourier transforms $X(\omega)$ and $H(\omega)$ using the following reasoning. \begin{align*} x(t) &= t\, e^{-2t} \,u(t)\\ & = t\; y(t) \qquad \text{where $y(t) := e^{-2t} \,u(t)$}\\ \implies X(\omega)&= i\, \frac{\partial Y}{\partial \omega}\qquad \text{where $Y$ is the F.T. of $y(t)$}\\ &= \frac{1}{(2+ i \omega)^2}. \end{align*} $H(\omega)$ is simply $1/(4+i\omega)$. Now \begin{align*} (x*h)(t) \stackrel{\text{FT}}{\longleftrightarrow} X(\omega) H(\omega) = \frac{1}{(2+i\omega)^2(4+i\omega)}. \end{align*}
If at this point I were to employ the inverse Fourier transform, I would obtain $(x*h)(t) = \frac{1}{2\pi} \int_{-\infty}^\infty \frac{e^{i \omega t}}{(2+i\omega)^2(4+i\omega)} \, \mathrm d\omega$. But this seems quite a complicated integral (considering this was a 4 mark question in an exam). Is there a simpler way to go about obtaining the inverse Fourier transform? Am I going about the question correctly?
I appreciate any help.
• My impulse would be to evaluate that last integral by turning it into a complex contour integral. Since the integral only has two poles above the real line, this shouldn’t be hard. – Semiclassical Jan 19 '18 at 14:14
• Also, how are you defining the Fourier transform? In the version I know, the Fourier transform of $(f*g)(t)$ would be $F(\omega)G(\omega)$ not $F(i\omega)G(i\omega)$. – Semiclassical Jan 19 '18 at 16:14
• @Semiclassical sorry about that, it's a notation purely to show that the argument may be complex also, equivalently it's $X(\omega)$. – Luke Collins Jan 19 '18 at 18:17
I think you're expected to apply the partial fraction decomposition: $$\frac 1 {(4+i \omega)(2+i \omega)^2} = \frac 1 {4(4+i \omega)} - \frac 1 {4(2+i \omega)} + \frac 1 {2(2+i \omega)^2},$$ and you already have the functions whose transforms produce those terms. | 2019-07-18T13:36:56 | {
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https://mathematica.stackexchange.com/questions/252994/does-this-series-converge-or-diverge | # Does this series converge or diverge?
Let us consider the series $$\sum_{n=1}^\infty \frac {\{\sqrt n\}-\frac 1 2} n$$, where $$\{\cdot\}$$ denotes the fractional part of a real number. On the one hand,
SumConvergence[(FractionalPart[Sqrt[n]] - 1/2)/n, n]
False
On the other hand,
NSum[(FractionalPart[Sqrt[n]] - 1/2)/n, {n, 1, Infinity}]
-0.725204
Let us investigate it. The output of
N[Table[Sum[(FractionalPart[Sqrt[n]] - 1/2)/n, {n, 1, 10^k}], {k, 1, 5}], 15]
{-0.679200513405746, -0.755781168886997, -0.803827960877609, -0.810978369858450, -0.816045768107535}
suggests the convergence, but the sum of the series does not equal -0.725204. Let us look at the plots
DiscretePlot[(FractionalPart[Sqrt[n]] - 1/2), {n, 1, 121}]
and
DiscretePlot[(FractionalPart[Sqrt[n]] - 1/2), {n, 121, 144}]
The plots show that approximately from k^2 to (k^2+(k+1)^2)/2 (k is a positive integer) the terms are negative and from (k^2+(k+1)^2)/2 to (k+1)^2 the terms are nonnegative. Let us estimate the sums over these intervals by
AsymptoticSum[1/n, {n, (k^2 + (k + 1)^2)/2, (k + 1)^2}, k -> Infinity]
1/k
and
AsymptoticSum[1/n, {n, k^2, (k^2 + (k + 1)^2)/2}, k -> Infinity]
1/k
This also suggests the convergence. However, the above is a plausible reasoning, not a proof.
The questions are: how to accurately prove or disprove the convergence with Mathematica? how to numerically calculate the sum of the series under consideration?
• A harder question is about the series $$\sum_{n=1}^\infty \frac {\{\sqrt n\}-\frac 1 2} ne^{2\pi i t \log n} ,$$ where $t\in\mathbb R$ . Aug 6 '21 at 18:28
• BTW, SumConvergence[FractionalPart[Sqrt[n]]/n, n] results in True instead of False. Aug 7 '21 at 12:21
Here is my answer to my question done with Mathematica. If n>=k^2 && n<(k+1)^2, where k is a positive integer, then Floor[Sqrt[n]]==k and FractionalPart[Sqrt[n]]==Sqrt[n]-k. Now we consider
Sum[(Sqrt[n] -k-1/2)/n, {n,k^2,(k + 1)^2 - 1},Assumptions -> k\[Element] PositiveIntegers]
1/2 (2 HurwitzZeta[1/2, k^2] - 2 HurwitzZeta[1/2, (1 + k)^2] + PolyGamma[0, k^2] + 2 k PolyGamma[0, k^2] - PolyGamma[0, 1 + 2 k + k^2] - 2 k PolyGamma[0, 1 + 2 k + k^2])
Series[%, {k, Infinity, 2}]
-(2/(3 k^2))+O[1/k]^3
This implies the series under consideration converges. In order to obtain the numerical value of the sum, we find
NSum[1/2 (2 HurwitzZeta[1/2, k^2] - 2 HurwitzZeta[1/2, (1 + k)^2] +
PolyGamma[0, k^2] + 2 k PolyGamma[0, k^2] -
PolyGamma[0, 1 + 2 k + k^2] -
2 k PolyGamma[0, 1 + 2 k + k^2]), {k, 1, Infinity}]
-0.817595
• I'd like to elaborate the proof. Because of Series[Sum[(Sqrt[n] - k - 1/2)/n, {n, k^2, (k + 1/2)^2}], k -> Infinity] which results in -(1/(4 k))+O[1/k]^2 and Series[Sum[(Sqrt[n] - k - 1/2)/n, {n, (k + 1/2)^2, (k + 1)^2 - 1}], k -> Infinity] which results in 1/(4 k)+O[1/k]^2, the partial sum $S_n$ for $n\in [k^2,(k+1)^2]$ satisfies the estimates $$S_n \ge S_{k^2}- 1/(4k)+O(1/k^2),\,S_n \le S_{k^2}+1/(4k)+O(1/k^2).$$ The squeeze theorem ends the proof. Aug 8 '21 at 4:32 | 2022-01-22T08:13:50 | {
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http://mathhelpforum.com/math-puzzles/133126-grains-rice-chess-board.html | # Math Help - Grains of rice on a chess board
1. ## Grains of rice on a chess board
Hi, firstly apologies if this is not the correct section.
At work today I was told that if you took a chess board (with 64 squares) and you were to place 1 grain of rice on the 1st square, 2 on the second square, 4 on the 3rd square, 8 on the 4th and so on doubling the grains of rice on each successive square compared to the previous square - that there would not be enough grains of rice in the world to do this.
Whilst im skeptical that there would not be enough rice in the world to do this...
how many grains of rice would be on the 64th square?, and
what would the formula be for working this out?
I guess I could work this out just by doubling 1 then 2 then 3 etc, but Id rather learn something new, so if there is a formula could someone explain it aswell.
As you can probably guess Im not the sharpest when it comes to maths, so thanks for your time and help.
Chez
2. Welcome to MHF.
If you call "n" the nth square, there are $2^{n-1}$ grains in the nth's square.
So in the last square there is $2^{63}$ grains. To see how big or how small this number is, go there: http://www.wolframalpha.com/input/?i=2^63.
I don't really know the number of rice's grains in the world, though it shouldn't be so hard calculating it. If you can give us the rice world production in tons by year, we will have the order of magnitude of the number of rice grains in the world, assuming that 1 kilo contains a certain number of grains.
Notice also that on the 63 th square there is half of the grains in the 64th square and that you also have to sum it up.
This number (of rice on the last square), as huge as it may look like, doesn't scar me. Indeed, there is 100 000 times more atoms of lead in a mole of lead, namely a small fragment of lead that you could have in one hand.
3. Originally Posted by arbolis
Welcome to MHF.
If you call "n" the nth square, there are $2^{n-1}$ grains in the nth's square.
So in the last square there is $2^{63}$ grains. To see how big or how small this number is, go there: http://www.wolframalpha.com/input/?i=2^63.
I don't really know the number of rice's grains in the world, though it shouldn't be so hard calculating it. If you can give us the rice world production in tons by year, we will have the order of magnitude of the number of rice grains in the world, assuming that 1 kilo contains a certain number of grains.
Notice also that on the 63 th square there is half of the grains in the 64th square and that you also have to sum it up.
This number (of rice on the last square), as huge as it may look like, doesn't scar me. Indeed, there is 100 000 times more atoms of lead in a mole of lead, namely a small fragment of lead that you could have in one hand.
Thanks that answers my question nicely.
World rice production is around 600 million tons per year (presumably this is a metric ton).
There are approx 36 950 grains of rice in a kilogram
so 36 950 000 grains in a ton,
so grains of rice produced in the world in a year is 36 950 000 * 600 = 22170000000 so something short of whats needed.
Unless my maths is wrong (likely) so please correct if thats the case.
Thanks again
4. Originally Posted by arbolis
Welcome to MHF.
If you call "n" the nth square, there are $2^{n-1}$ grains in the nth's square.
So in the last square there is $2^{63}$ grains. To see how big or how small this number is, go there: http://www.wolframalpha.com/input/?i=2^63.
I don't really know the number of rice's grains in the world, though it shouldn't be so hard calculating it. If you can give us the rice world production in tons by year, we will have the order of magnitude of the number of rice grains in the world, assuming that 1 kilo contains a certain number of grains.
Notice also that on the 63 th square there is half of the grains in the 64th square and that you also have to sum it up.
This number (of rice on the last square), as huge as it may look like, doesn't scar me. Indeed, there is 100 000 times more atoms of lead in a mole of lead, namely a small fragment of lead that you could have in one hand.
Yes, but the number of grains of rice in total would be $\sum_{n=0}^{64}2^n=2^65-1$
5. Originally Posted by Drexel28
Yes, but the number of grains of rice in total would be $\sum_{n=0}^{64}2^n=2^65-1$
Yes, that's why I wrote "you have to sum it up".
There's a slight error, if you start at n=0, you must end at n=63 since there are 64 squares but you are right we must sum all the rices on every square.
6. Originally Posted by Drexel28
Yes, but the number of grains of rice in total would be $\sum_{n=0}^{64}2^n=2^65-1$
The volume ocupied by the grains on the last square alone would fill a sphere of radius the order of $3.6$ km (or more if a packing fraction less than 1 is used)
CB
7. Originally Posted by chezza
Thanks that answers my question nicely.
World rice production is around 600 million tons per year (presumably this is a metric ton).
There are approx 36 950 grains of rice in a kilogram
so 36 950 000 grains in a ton,
so grains of rice produced in the world in a year is 36 950 000 * 600 = 22170000000 so something short of whats needed.
Unless my maths is wrong (likely) so please correct if thats the case.
Thanks again
I think you made an error in the expression "36 950 000 * 600 = 22170000000". You're multiplying the number of rice grains contained in one ton by 600 instead of 600 millions. Thus the final result should be 10^6 times greater than your result. So yes, it seems that with one year world production, you would have approximately 1/1000th of the grains needed to fill the last square, if I didn't make any arithmetic mental error.
8. Thanks for the replies,
Its been interesting reading though I have no idea what the formula given by drexel means
9. Originally Posted by chezza
Thanks for the replies,
Its been interesting reading though I have no idea what the formula given by drexel means
$\sum_{n=0}^{63}2^n=2^0+2^1+2^2+2^3+...+2^{62}+2^{6 3}=1+2+4+8+...+9223372036854775808$ which is the total number of rice grains. Each term of the sum represent the number of rice grain(s) in each square. Summing them up reach the number of rice grains in the whole board. | 2015-02-27T23:39:07 | {
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http://math.stackexchange.com/questions/35136/does-fx-in-mathbbzx-have-the-same-roots-as-fx-in-mathbbf-px | # Does $f(x) \in \mathbb{Z}[x] \$ have the same roots as $f(x) \in \mathbb{F}_p[x] \quad$?
Does $f(x)\in\mathbb{Z}[x]\$ have the same roots of $f(x)\in\mathbb{F}_p[x] \quad$?
Do the roots of a polynomial change when the coefficients of the polynomial are considered as elements of one ring or another? Would the answer to this change if the rings contained each other, as in the case of $f(x) \in \mathbb{Z}[x] \ \$ vs. $\ f(x) \in \mathbb{Q}[x] \quad$?
-
Short answer: Yes, the set of roots can change. $f(x)=x^2+1$ has no roots in $\mathbb{Z}$, but (the image of $f(x)$ under reduction modulo $p$) has roots in $\mathbb{F}_p$ if $p=2$ or $p\equiv 1 \pmod{4}$.
Longer answer: First, you have to be careful. If $f(x)$ is a polynomial with integer coefficients, then it is not literally a polynomial with coefficients in $\mathbb{F}_p$. Rather, the natural map "reduction modulo $p$", which is a ring map from $\mathbb{Z}$ to $\mathbb{F}_p$, induces a map from $\mathbb{Z}[x]$ to $\mathbb{F}_p[x]$, and you can consider the image of $f(x)$ in $\mathbb{F}_p[x]$.
Under this map, if $a\in\mathbb{Z}$ is a root of $f(x)$, then $a\bmod p$ is a root of $f(x)\bmod p$ in $\mathbb{F}_p[x]$. But the converse certainly does not hold; for one thing, there are infinitely many integers that are preimages of $a\bmod p$. And for another, no: you can have $f(x)\bmod p$ have roots in $\mathbb{F}_p[x]$, but no roots in $\mathbb{Z}[x]$, or only some roots. For example, for odd primes $p$, $f(x) = x^{p}-x$ has only two roots in $\mathbb{Z}$ ($x=0$ and $x=1$), but (its image under reduction modulo $p$) has $p$ roots in $\mathbb{F}_p$. So roots in $\mathbb{Z}$ yield roots in $\mathbb{F}_p$ (via reduction modulo $p$), but not conversely. Another example: $x^2+1$ has no roots in $\mathbb{Z}$, but (its reduction modulo $p$) has roots in $\mathbb{F}_p$ for every $p$ that is not congruent to $3$ modulo $4$.
More generally, the universal property of the polynomial ring gives easily:
Theorem. Let $R$ and $S$ be commutative rings. If $h\colon R\to S$ is a ring homomorphism, then $h$ induces a ring homomorphism $\overline{h}\colon R[x]\to S[x]$ of polynomial rings by $$\overline{h}(r_0+r_1x + \cdots + r_nx^n) = h(r_0) + h(r_1)x + \cdots + h(r_n)x^n.$$ Moreover, the map $\overline{h}$ commutes with evaluation maps in the following sense: if $\mathrm{eval}_a\colon R[x]\to R$ is the map $\mathrm{eval}_a(f(x)) = f(a)$, and $\mathrm{eval}_{h(a)}\colon S[x]\to S$ is the evaluation map at $h(a)$, then we have a commutative diagram $$\begin{array}{rcl} R[x] & \stackrel{\overline{h}}{\longmapsto} & S[x]\\ &&\\ \mathrm{eval}_a\downarrow & & \downarrow \mathrm{eval}_{h(a)}\\ R &\stackrel{h}{\longmapsto} & S\end{array}$$ for all $a\in R$.
In particular, if given $f\in R[x]$ and $g\in S[x]$ we let $$\mathrm{Roots}_R(f) = \{r\in R\mid f(r)=0\}$$ and $$\mathrm{Roots}_S(g) = \{s \in S\mid g(s)=0\},$$ then $$h\left(\mathrm{Roots}_R(f)\right) \subseteq \mathrm{Roots}_S(\overline{h}(f)),$$ but the inclusion may be proper.
For examples with proper inclusion in the setting of reduction modulo $p$, consider $f(x)=x^2+1\in\mathbb{Z}[x]$ which has not roots, but whose reduction modulo $p$, with $p=2$ or $p\equiv 1\pmod{4}$ has roots. Or $f(x)=x^p-x$, $p$ an odd prime, which has two roots in $\mathbb{Z}$ but its reduction modulo $p$ has $p$ roots in $\mathbb{F}_p$ (by Fermat's Little Theorem). For examples of proper inclusion with fields and $h$ an inclusion map, take $x^2+1\in\mathbb{R}[x]$, and consider its image in $\mathbb{C}[x]$; $x^2+1$ has no real roots, but its image in $\mathbb{C}[x]$ has roots.
Added. When dealing with inclusion, it usually makes more sense to consider the set of roots in the larger field, and just ask which ones lie in the smaller field. So we think of $x^2+1$ as a polynomial with complex coefficients, and then ask what $\mathrm{Roots}_{\mathbb{C}}(x^2+1)\cap\mathbb{R}$ is (this will equal $\mathrm{Roots}_{\mathbb{R}}(x^2+1)$).
Also, if $h(a)$, the image of $a$ is a root of $\overline{h}(f)$, this does not imply, in and of itself, that $a$ is a root of $f(x)$ (though it does if $h$ is one-to-one).
To summarize the longer answer: if $h\colon R\to S$ is a ring homomorphism, then the image of any root of $f(x)$ is a root of the image of $f(x)$, but there may be roots of the image that are not images of roots; and not everything in $R$ whose image is a root of the image of $f(x)$ is a root of $f(x)$ in general. But if $h$ is one-to-one, then $h(a)$ is a root of $\overline{h}(f(x))$ if and only if $a$ is a root of $f(x)$.
-
Holy crap, you wrote a book for me. Thanks, I appreciate it :) – badatmath Apr 26 '11 at 4:29
Hello, and thanks :) I guess I am confused because when you talk about polynomials like $x^2+1$ with coefficients in $\mathbb{Z}[x]$, you say this polynomial has roots, but its roots are in $\mathbb{C}$. If you say $x^2+1$ has coefficients in $\mathbb{F}_p[x]$, can you also say this polynomial has roots in $\mathbb{C}$? Does it make a difference that $\mathbb{C}$ does not contain $\mathbb{F}_p$? – badatmath Apr 26 '11 at 4:32
@badatmath: I did not say "the roots are in $\mathbb{C}$." The inclusion $\mathbb{Z}\hookrightarrow\mathbb{C}$ means you can consider polynomials with integer coefficients as polynomials with complex coefficients. If you do that, then the image of $x^2+1$ has roots when you consider it as an element of $\mathbb{C}[x]$, but no roots when you consider it as an element of $\mathbb{Z}[x]$. No, you cannot say that a polynomial in $\mathbb{F}_p[x]$ has roots in $\mathbb{C}$; that makes no sense precisely because $\mathbb{C}$ does not contain $\mathbb{F}_p$. – Arturo Magidin Apr 26 '11 at 4:35
@badatmath: But the reason $x^2+1$ has roots in $\mathbb{F}_p$ when $p\equiv 1 \pmod{4}$ is that in those situations, $-1$ is a square modulo $p$. For example, in $\mathbb{F}_5$, but $2\bmod 5$ and $3\bmod 5$ are roots of (the image of) $x^2+1$. – Arturo Magidin Apr 26 '11 at 4:36
Oh, okay. So you can't say that because the inclusion map does not exist. Thanks. – badatmath Apr 26 '11 at 5:18
Last question first, $2x-1$ certainly doesn't have the same roots in $\bf Z$ as in $\bf Q$. I'd say $x-2$ has the same root in $\bf Z$ as in $\bf Q$, namely, $2$, but I can imagine some situation in which one might want to distinguish between $2$ as an integer and $2$ as a rational.
Does $x-7$ have the same root in $\bf Z$ as in ${\bf F}_5$? I think it would be an abuse of notation to say yes. But abuses of notation can be very useful, so I wouldn't rule it out; I'd just make absolutely sure to put all my cards on the table first, that is, to say something like we identify elements of $\bf Z$ with their images in ${\bf F}_5$'' first.
-
No. In fact, this failure can be rather extreme, as for polynomial rings over finite fields there are nonzero polynomials with integer coefficients that vanish at every point in the field. Easy examples are provided by encoding boolean contradictions as polynomials in $\mathbb{F}_2[X]$, such as $(x \text{ and } x)\text{ xor } x$ which is represented as $x^2 + x$ or $((x \text{ and } x)\text{ and } x) \text{ xor } x$ which is represented as $x^3 + x$. You can also get such polynomials in $\mathbb{F}_p[X]$ by taking the product $(x - (p - 1))(x - (p - 2))\cdots (x - 1)x$.
-
Echoing Arturo, since there's a surjection $\phi: \mathbb{Z}[x] \to \mathbb{F}_p[x]$, if $a$ is a root of $f$ in $\mathbb{Z}[x]$, then $\phi(a)$ is a root of $f$ in $\mathbb{F}_p[x]$. So a root in $\mathbb{Z}[x]$ translates into a root in $\mathbb{F}_p$, but they can overlap. For example, $x^2 - 1 = (x+1)(x-1)$ has two distinct roots $1,-1$ in $\mathbb{Z}[x]$ but one repeated root $1$ in $\mathbb{F}_2[x]$.
In the same line as Alex's comment, you have to be a little careful. If $p$ divides every coefficient of $f$, then $f \equiv 0$ in $\mathbb{F}_p[x]$.
In addition, some irreducible elements in $\mathbb{Z}[x]$ do factor in $\mathbb{F}_p[x]$. See Arturo's answer.
In general, it's easier to factor in $\mathbb{F}_p[x]$ because there are fewer elements to check divisibility by. And that can help determine factorization in $\mathbb{Z}[x]$. A theorem states that if the leading coefficient of $f$ is not divisible by $p$, and if $f$ is irreducible in $\mathbb{F}_p$, then $f$ is irreducible in $\mathbb{Z}$ (Artin 2nd Edition 12.4.3). The converse does not hold for the same reasons as discussed in the previous paragraph.
Factorization in $\mathbb{Z}[x]$ and $\mathbb{Q}[x]$ is equivalent up to constant multiples, but linear factors are always roots in $\mathbb{Q}[x]$, while they aren't necessary so in $\mathbb{Z}[x]$, as Gerry stated. Roots in $\mathbb{Z}[x]$ will be roots in $\mathbb{Q}[x]$, but not necessarily the other way around. | 2015-10-05T04:41:39 | {
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https://math.stackexchange.com/questions/3112037/what-is-the-largest-integer-value-of-n-for-which-8n-evenly-divides-100 | # What is the largest integer value of $n$ for which $8^n$ evenly divides $(100!)$?
I know that this may be an unnecessary question, but I am a bit confused. The problem asks for the highest integer $$n$$ such that $$8$$ to the power of $$n$$ is divisible, evenly of course, by $$100$$. Now, I searched the site, and, in general, it seems that one can use floor function for a problem like this, but this seems to only work for prime numbers possibly. My process, which I realized was incorrect:
The floor function of $$100/8 = 12$$, and then doing it for the second power would lead to one, and, by adding those up, I acquired thirteen. Of course, after seeing the answer, $$32$$, I went back to see what was wrong and did the problem slower. I got $$12$$ numbers from the numbers in $$100!$$, and then got another $$8$$ from $$2 \times 4$$, but, that can be applied for all the multiples of $$2$$ and $$4$$ that aren't of $$8$$. So, essentially, I am wondering if there is a quicker method for calculating this number without specifically counting out the numbers. Thanks in advance!
• are you talking about (100!) or (100) ?
– Seth
Feb 13 '19 at 22:59
• I intended for it to be meant as (100!) with a question mark at the end. Feb 13 '19 at 23:00
• Edited title for easier reading. Thanks for the concern. Feb 13 '19 at 23:01
• Is it $8^n$ is divisible by $100$ or $8^n$ divides $100!\,$? Feb 13 '19 at 23:07
• It is the second one ( 8^n divides (100!) ) Feb 13 '19 at 23:09
I think the easiest way to answer this question is to factorize $$100!$$. Actually, a partial factorization will be sufficient. Thus, we see that $$100! = 2^{97} \times 3^{48} \times 5^{24} \times \ldots \times 83 \times 89 \times 97.$$
Since $$8 = 2^3$$, we need to divide $$97$$ by $$3$$ and discard the remainder. That is, rewrite $$2^{96}$$ as $$8^n$$ and there's your answer.
• Thanks for the answer! I was hoping that there was a quicker method (maybe an algorithmic one). Feb 13 '19 at 23:14
• Thanks for the suggestion (I'm new to the site, so, I don't know how to make this work). Feb 13 '19 at 23:15
• This suggests to me the following algorithm: repeatedly divide the factorial by 2, increment a counter as you do so, until you get an odd number, specifically a number from oeis.org/A049606 Feb 14 '19 at 2:44
It's easiest, I think, to do this with powers of $$2:$$
$$\left\lfloor{100\over2}\right\rfloor+ \left\lfloor{100\over4}\right\rfloor+ \left\lfloor{100\over8}\right\rfloor+ \left\lfloor{100\over16}\right\rfloor+ \left\lfloor{100\over32}\right\rfloor+ \left\lfloor{100\over64}\right\rfloor=97=32\cdot3+1$$
so the greatest exponent of $$8$$ is $$32$$.
• Thanks for the answer! Feb 13 '19 at 23:13
• This is known as Legendre's formula. Feb 13 '19 at 23:35
• @Bernard Reading mathworld.wolfram.com/LegendresFormula.html I'm not quite sure this is right, though it's definitely related. Of course it's also possible I have misunderstood. Feb 14 '19 at 23:58
Just sum up the $$2$$-adic orders of the even numbers from $$2$$ to $$100$$
You get that $$(100)! = 2^{97}\cdot A=8^{32}\cdot 2A$$ where $$A$$ is the product of all the remaining (odd) factors of $$100!$$
• Thanks for taking the time to answer! Feb 13 '19 at 23:20
• You're welcome. Hope it helps. Feb 13 '19 at 23:21
So you're looking for roughly a third of $$n = \sum_{i = 1}^\infty \left\lfloor \frac{m}{p^i} \right\rfloor,$$ where $$m = 100$$ and $$p$$ is 2 ($$m$$ can be any positive integer and $$p$$ can be any odd prime).
Of course you don't actually have to even try to go to infinity. As soon as you notice $$\frac{m}{p^i} < 1,$$ you can stop the iteration.
The important thing to understand here is that each consecutive number by which you multiply to get a factorial "adds" its prime factors' exponents to the factorial's prime factors' exponents.
For example, $$11! = 2^8 \times 3^4 \times 5^2 \times 7 \times 11 \times 13^0 \times 17^0 \times 19^0 \times \ldots$$ Since $$12 = 2^2 \times 3^{(1)}$$, we can just add 2 to 2's exponent and 1 to 3's exponent, to obtain $$12! = 2^{10} \times 3^5 \times 5^2 \times 7 \times 11 \times 13^0 \times$$ $$17^0 \times 19^0 \times \ldots$$ And then the factorization of 13! is a simple matter of adding 1 to 13's exponent.
So in the specific case you're looking at, the first question is: how many numbers less than or equal to 100 are even? Half of them, giving you 50. But this doesn't account for how many of them are "doubly even," since they contribute at least 2 each to 2's exponent in the factorization of 100! So there's 25 of those, bringing our tally up to 75.
And then twelve of them are divisible by 8 (tally's 87 now), six are divisible by 16 (tally 93), three are divisible by 32 (tally 96), only one one is divisible by 64 (tally 97), and none are divisible by 128, 256, 512, 1024, etc. (tally stays at 97).
All that's left to do now is to solve $$2^{97} = 8^x$$, and round down $$x$$ if necessary.
Well $$8^n = 2^{3n}$$ so it will do to find the highest power of $$2$$ which divide $$100!$$.
Which means considering the prime factorization of $$100!= 2^k*3^j*5^m.....$$. What is $$k$$ in the $$2^k$$?
Well $$100! = 1 * 2 * 3* 4* ....... *97*98*99*100$$ but only the even numbers contribute to powers of $$2$$.
So $$100! = 2*4*6*......*98*100*($$ a bunch of odd terms$$)=$$
$$= (2*1)*(2*2)*(2*3)*..... *(2*49)*(2*50)*($$ the odd terms $$)=$$
$$= 2^{50}*(1*2*3*4*.....*50)*($$ the odd terms $$)$$.
Now the even terms of $$1,2,3,4,...., 50$$ are going to contribute to powers of $$2$$ so
$$100! = 2^{50}*(2*4*6*...*48*50)*($$ the odd terms up to fifty$$)*($$ a heck of a lot of odd terms$$)=$$
$$2^{50}*(2*1)*(2*2)*(2*3)*.....*(2*24)*(2*25)*($$ just a huge honkin amount of odd terms that we hae utterly no need to keep track of$$)=$$
$$2^{50}*2^{25}*(1*2*3.....*25)*($$ sh!tload of odd terms$$)=$$
$$2^{50}*2^{25}*(2*4*....*24)*($$ something odd $$)=$$
$$2^{50}*2^{25}*2^{12}*(1*2*3*....*12)*($$ odd monster $$)=$$
$$2^{50}*2^{25}*2^{12}*(2*4*6*8*10*12)*($$ odd thing $$)=$$
$$2^{50}*2^{25}*2^{12}*2^6*(1*2*3*4*5*6)*$$ ODD$$)=$$
$$2^{50}*2^{25}*2^{12}*2^6*(2*4*6)*$$ MEGA-ODD$$)=$$
$$2^{50}*2^{25}*2^{12}*2^6*2^3*(1*2*3)*$$ ODDasaurus $$)=$$
$$2^{50}*2^{25}*2^{12}*2^6*2^3*2^1*$$ MEGA-MECHA-ODD-atron-a-galactadingy.
So $$100! = 2^{50+25+12+6+3+1}*M$$ where $$M$$ is an odd number.
$$100! = 2^{97}*M= 2^{3*32}*2*M= 8^{32}*2*M$$.
So $$8^{32}|100!$$ but $$8^{33}$$ does not.
====
In general.
To find the highest power of prime $$p$$ that divides $$N!$$ realize that of the terms $$1..... N$$ that $$[\frac Np]$$ of those terms are mulitples of $$p$$ so $$p^{[\frac Np]}$$ will divide $$N!$$.
However $$[\frac N{p^2}]$$ of those terms are not just multiples of $$p$$ but of $$p^2$$ and these term contribute $$[\frac N{p^2}]$$ additional powers of $$p$$.
Repeating until we are done:
The highest power of $$p$$ dividing $$100!$$ will be $$p^{[\frac Np] + [\frac N{p^2}] + [\frac N{p^3}] + ......}$$.
In the case of $$2$$ and $$100!$$ it was
$$2^{[\frac {100}2] + [\frac {100}4 ]+ [\frac {100}8] + [\frac {100}{16}[ + [\frac {100}{32}] + [\frac {100}{64}]}=2^{50 + 25 + 12+ 6 + 3+ 1}$$. | 2021-09-26T23:21:28 | {
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https://math.stackexchange.com/questions/2817987/finding-absolute-minimum-and-absolute-maximum-of-fx-y-xy | # Finding Absolute Minimum and Absolute Maximum of $f(x,y)=xy$
Let $\ f(x,y)=xy$. Use the method of Lagrange multipliers to find the maximum and minimum values of the function f on the circle $\ x^2+y^2=1$
First we note that the function $f$ is continuous and the set $S={(x,y):x^2+y^2=1}$ is compact, hence extrema are guaranteed. Using the method Lagrange multipliers, I set $\nabla f=\lambda\nabla g$, where $g(x,y)=x^2+y^2-1$. Following through the calculations, I arrived at four critical points: $$\Big(\frac{1}{\sqrt{2}},\frac{1}{\sqrt{2}}\Big),\Big(\frac{1}{\sqrt{2}},-\frac{1}{\sqrt{2}}\Big),\Big(-\frac{1}{\sqrt{2}},\frac{1}{\sqrt{2}}\Big),\Big(-\frac{1}{\sqrt{2}},-\frac{1}{\sqrt{2}}\Big)$$ Substituting these points into the function $f$, I obtained a maximum at $$\Big(\pm\frac{1}{\sqrt{2}},\pm\frac{1}{\sqrt{2}}\Big)=\frac{1}{2}$$ and a minimum at $$\Big(\pm\frac{1}{\sqrt{2}},\mp\frac{1}{\sqrt{2}}\Big)=-\frac{1}{2}$$
My question is, how do we now find the absolute maximum and absolute minimum of the function $f$ on the unit disc $x^2+y^2\leq 1$?
My attempt so far:
We want to find all the critical points. So to find stationary points, we set $$\nabla f=\vec{0}$$ Solving this, we find that $(0,0)$ is a stationary point. So, $f(0,0)=0$. Hence the absolute maximum is $\frac{1}{2}$ and the absolute minimum is $-\frac{1}{2}$, as these are all the critical points of $f$. This does not sit well with me, as I am unsure of my working/logic. Can this be improved on?
• Could be beneficial to parameterize the unit circle $(x,y)=(\cos t,\sin t)$ and consider optimizing $f(x,y)=xy=\cos t\sin t=\frac12\sin2t$ for $t\in[0,2\pi]$. – A.Γ. Jun 13 '18 at 8:10
• an observation is that under the constraint: $(x+y)^2 = x^2+y^2+2xy = 1 + 2xy$ and the solution becomes trivial. – Stefan Jun 13 '18 at 8:17
• You probably mean the unit disk $x^2+y^2\le 1$ ? – user65203 Jun 13 '18 at 8:18
• Yes, question has been edited. Thanks! – user557493 Jun 13 '18 at 8:21
The method of Lagrange multipliers tells us the the local extremes of the restriction of $f$ to circle are the ones that you got. So, the absolute maximum and the absolute minimum must be attained at some of them (the absolute maximum and the absolute minimum are local extremes). Since $f$ takes the value $\frac12$ in two of them and no value greater than $\frac12$, $\frac12$ is necessarily the maximum. The same argument applies to the minimum.
• Thank you for your response! Just so that I understand, are there 5 critical points? Two of which $f$ takes the value $\frac{1}{2}$, two of which $f$ takes the value $-\frac{1}{2}$ and a final critical point which takes the $f$ value of $0$? I am unsure of how many critical points there are. – user557493 Jun 13 '18 at 8:09
• No. There are only four critical points. None of the points at which $f$ takes the values $0$ (that is, the points $(\pm1,0)$ and $(0,\pm1)$) is critical. – José Carlos Santos Jun 13 '18 at 8:11
• The following question asks to classify every stationary point of the function $f$. Are you able to provide some insight into this question? Is $(0,0)$ a stationary point? How many stationary points are there? – user557493 Jun 13 '18 at 8:17
• Is $(0,0)$ a saddle point? Also, if $f(0,0)\geq\frac{1}{2}$ would this be the absolute max? I thought the absolute max of $f$ was the largest value of every critical point evaluated at $f$. – user557493 Jun 13 '18 at 8:24
• This is a different question. Please post it as such. I think that I've answered the original question. – José Carlos Santos Jun 13 '18 at 8:26
You have the minimization optimization problem of the function $f(x,y) = xy$ over the space $S= \{ (x,y) \in \mathbb R^2 : x^2 + y^2 = 1\}$. A known way to deal with Lagrange multipliers is by the Kuhn-Tucker Lagrange method.
First of all, observe that $f(x,y)$ is continuous and smooth and that the space $S$ is compact. Thus, this means that there exists a minimum $(\bar{x},\bar{y})$ for $f(x,y)$ in $S$.
By the Kuhn-Tucker Lagrange method, we yield :
$$f_0(x,y) = xy, \; \; f_1(x,y)= x^2+y^2-1$$
and then the K.T.L. system :
$$\begin{cases} \nabla f_0 + \lambda_1\nabla f_1 = 0 \\ \lambda_1 f_1 =0 \end{cases} \Rightarrow \begin{cases} \begin{bmatrix} y \\ x \end{bmatrix} + \lambda_1\begin{bmatrix} 2x \\ 2y \end{bmatrix} =0 \\ \lambda_1(x^2 + y^2 -1) \;= 0\end{cases}$$
Check cases for $\lambda_1 = 0$ and $\lambda_1 >0$ and then you'll yield the same results. (Maximum is given for applying the same method for $-f(x,y)$ or simply you yield the same points as you did.
Now, if there existed another minimum or maximum, it should satisfy the K.T.L. problem. Since no other point satisfies it, these are all the minimums and maximums. Observing that you have two possible minimum and maximum points (since the values are equal) for $f(x,y)$ over $S$, this means that you have a total maximum and minimum at both of the points each time.
The Lagrange method gives you the local extrema of the function constrained to the boundary *. The stationary points give you the local extrema, and you consider those inside the boundary.
Then the global extrema are achieved by the local extrema that yield the largest/smallest values.
*Alternatively, you could use a parametric equation of the boundary, let $x=\cos t,y=\sin t$, and find the local extrema of $\cos t\sin t$, which occur at $t=\dfrac{k\pi}4$, giving values $\pm\dfrac12$.
• Okay. So, the largest/smallest values of the local extrema (those on the boundary, stationary points etc) determine the global extrema? Hence if for the stationary point $(0,0)$, $f(0,0)\geq\frac{1}{2}$, then $(0,0)$ would be the global maximum? Is this what you mean? – user557493 Jun 13 '18 at 8:35
• @Bell: right. Try with $1-x^2-y^2$ or $xy-x^2-y^2$. – user65203 Jun 13 '18 at 8:37
• I will indeed. That was very helpful. Thank you! – user557493 Jun 13 '18 at 8:38 | 2021-07-30T12:55:33 | {
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https://www.physicsforums.com/threads/how-do-i-complete-the-square.769226/ | # How do I complete the square?
1. Sep 5, 2014
### shreddinglicks
1. The problem statement, all variables and given/known data
I identify the curve by finding a Cartesian equation for the curve.
2. Relevant equations
r = 3sin(θ)
3. The attempt at a solution
r^2 = 3r(y/r)
x^2+y^2 = 3y
x^2-3y+y^2 = 0
Now I'm stuck, I don't remember how to complete the square as I haven't done it in ages.
2. Sep 5, 2014
### Mentallic
Notice that the only x in the equation is a single x2 term so that's already a square. We just need to turn y2-3y into a complete square now.
The binomial expansion for the following square value is
$$(y-a)^2 = y^2-2ay+a^2$$
So we want to choose the value of 'a' such that we have turned the right hand side into the $y^2-3y$ expression.
So essentially since we want $y^2-3y=y^2-2ay$ then comparing coefficients of y clearly we want -3=-2a, so a=3/2.
But we also need the a2 term in there in order to turn it back into a perfect square. If we add a2 to both sides of the equation then we keep everything balanced and we get
$$x^2+y^2-3y+(3/2)^2=(3/2)^2$$
Can you take it from here?
3. Sep 5, 2014
### shreddinglicks
I see, so I have radius r = 3/2. How do I get the center?
4. Sep 5, 2014
### Mentallic
Have you completed the square on y? The general form for the circle is
$$(x-a)^2+(y-b)^2=r^2$$
And the centre is located at (a,b). Note that $x^2=(x-0)^2$
5. Sep 5, 2014
### shreddinglicks
That is interesting. My textbook has me making a table using radian angles and plugging into the equation to find values of r and plotting them on a graph.
6. Sep 6, 2014
### Mentallic
Well the thing about parametric curves is that you can't always find a nice cartesian equation to represent them. In this case it's simply a circle and so doing what you did is best to find what the parametric curve actually is.
Also, it would be good practise to do what your textbook is expecting so that you can find the rough shape of curves that normally wouldn't be as nice.
7. Sep 6, 2014
### HallsofIvy
Staff Emeritus
you have $-3y+ y^2$. You should recall, or be able to calculate, that $(y- a)^2= y^2- 2ay+ a^2$. Comparing that with $y^2- 3y$, you see that your "2a" is 3. Since 2a= 3, a= 3/2 and $a^2= 9/4$. You need to add (and subtract so you won't actually change the value of the expression) 9/4: $y^2- 3y= y^2- 3y+ 9/4- 9/4= (y- 3/2)^2- 9/4$. | 2017-10-20T07:37:21 | {
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https://ferienwohnung-gluecksburg.net/rbtm7b/912670-how-to-find-the-area | The area of the square is . Area between Curves Calculator. In a similar way, you can find the area of any closed Polyline geometry with AREA command. Some teachers require students to use z-tables, but there are … Area of a Circle Calculator. See More. Area = ½ bh = ½ × 20 × 12 = 120 Knowing Three Sides There's also a formula to find the area of any triangle when we know the lengths of all three of its sides. Square Meter vs Meter Square You can find the area in square units of the rhombus by multiplying the lengths of the two diagonals (d 1 and d 2) and dividing by two. For FREE. You'll need: Tape measurer or ruler Calculator You may also be interested in: How to Calculate the Volume of a Cube - Formula and Examples. We will use the equation in the picture below. 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Is irregular, but we can break it into two rectangles the most formulas... Square area of a circle, however, or planar lamina, in picture... ( sq rectangle has an area of a circle to find the area for each shape.. | 2021-10-19T17:52:33 | {
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http://slideplayer.com/slide/4261794/ | # If f(x) 0 for a x b, then b f(x) dx a is defined to be the area of the region under the graph of f(x) on the interval [a,b]. The precise definition.
## Presentation on theme: "If f(x) 0 for a x b, then b f(x) dx a is defined to be the area of the region under the graph of f(x) on the interval [a,b]. The precise definition."— Presentation transcript:
If f(x) 0 for a x b, then b f(x) dx a is defined to be the area of the region under the graph of f(x) on the interval [a,b]. The precise definition of this integral as the area is typically done by defining the integral as a limit of sums; this leads to the Fundamental Theorem of Calculus, which states that such integrals can be calculated using anti-differentiation. The definition of this integral can then be extended so that f(x) may possibly take on negative values (and areas below the x axis will be negative).
If f(x,y) 0 for all (x,y) in a region D of the xy plane, then the double integral is defined to be the volume under the graph of f(x,y) on the region D. When the region D is a rectangle R = [a,b] [c,d], then the volume is that of the following solid: (b,c) z (a,c) (a,d) (b,d) z = f(x,y) y x f(x,y) dx dy= f(x,y) dA DD
ExampleSuppose f(x,y) = k where k is a positive constant, and let R be the rectangle [a,b] [c,d]. Find (b,c) z (a,c) (a,d) (b,d) f(x,y) = k Since the integral is the volume of a rectangular box with dimensions b – a, d – c, and k, then the integral must be equal to y x k(b – a)(d – c). f(x,y) dA. R ExampleLet R be the rectangle [0,1] [0,1]. Find (1,1,0) x y z Since the integral is the volume of half of a unit cube, then the integral must be equal to f(x,y) = (1 – x) (1,0,0) (0,0,1)(0,1,1) (0,1,0) 1/2. (1 – x) dA. R
Cavalieri’s Principle : Suppose a solid is sliced by a series of planes parallel to the yz plane, each labeled P x, and the solid lies completely between P a and P b. z y x A(x)A(x) PxPx Let A(x) be the area of the slice cut by P x. If x represents the distance between any pair of parallel planes, then A(x) x is approximately the volume of Summing the approximations A(x) x is an approximation to the portion of the solid between P x and its successor. the volume of the solid. Taking the limit of the sum as x 0, we find that the solid’s volume is b A(x) dx a
Apply Cavalieri’s principle to find the volume of the following cone: length = 12 length = 13 circle of radius z x y 2.5 y = 5 — x 12 For 0 x 12, the slice of the cone at x parallel to the yz plane is a circle with diameter The area of the slice of the cone at x is A(x) = 5 — x. 12 25 —– x 2. 576 The volume of the cone is A(x) dx = a b 0 12 25 —– x 2 dx = 576 25
To apply Cavalieri’s principle to the double integral of f(x,y) over the rectangle R = [a,b] [c,d], we first find, for each value of x, the area A(x) of a slice of the solid formed by a plane parallel to the yz plane: (b,c) x y z (a,c) (a,d) (b,d) z = f(x,y) A(x) = Consequently, d f(x,y) dy c f(x,y) dA = R b A(x) dx = a b d f(x,y) dydx. ac By reversing the roles of x and y, we may write f(x,y) dA = R d b f(x,y) dxdy. ca
These equations turn out to be valid even when f takes on negative values. The integrals on the right side of each equation are called iterated integrals (and sometimes the brackets are deleted). ExampleLet R be the rectangle [–1,1] [0,1]. Find (x 2 + 4y) dA. R (x 2 + 4y) dA = R (x 2 + 4y) dy dx = 0 11 –1–1 x 2 y + 2y 2 dx = y = 0 1 1 –1–1 (x 2 + 2) dx = – 1 1 x 3 /3 + 2x = x = –1 1 14 — 3 Alternatively, one could find that (x 2 + 4y) dx dy = 0 11 –1–1 14 — 3
ExampleLet R be the rectangle [0, /2] [0, /2]. Find (sin x)(cos y) dA. R (sin x)(cos y) dA = R (sin x)(cos y) dy dx = 0 /2 0 (sin x)(sin y) dx = /2 0 y = 0 /2 (sin x) dx = 0 /2 – cos x = x = 0 /2 1 Alternatively, one could find that (sin x)(cos y) dx dy = 0 /2 0 1
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Similar presentations | 2018-04-27T07:57:53 | {
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https://discourse.julialang.org/t/difference-between-f-1-0e-7-f-1-0e-8/32779 | # Difference between f(1.0e-7) & f(1.0e-8)
the limit of the following function at 0 is 1,
yet f(1.0e-8) gives 0 while f(1.0e-7) gives the correct answer, why?
julia> f(x)=2*(1-x*exp(-x)-exp(-x))/(1-exp(-x))^2
f (generic function with 1 method)
julia> f(1.0e-7)
0.9992008230547272
julia> f(1.0e-8)
0.0
The problem is catastrophic cancellation due to limited accuracy of double precision floating point arithmetic:
julia> 1 - (1e-7) * exp(-1e-7)
0.99999990000001
julia> exp(-1e-7)
0.999999900000005
julia> 1 - (1e-7) * exp(-1e-7) - exp(-1e-7)
4.9960036108132044e-15
julia> 1 - (1e-8) * exp(-1e-8)
0.9999999900000001
julia> exp(-1e-8)
0.9999999900000001
julia> 1 - (1e-8) * exp(-1e-8) - exp(-1e-8)
0.0
7 Likes
Note also that your f is probably quite wrong also for x = 1e-7:
julia> f(big"1e-7")
1.000000033333333333333322222222222222226190476190476189153439153575262308010033
Using textbook formulas involving very small numbers is often wrong in numerical computing.
4 Likes
The following implementation computes your g(\delta) accurately in double precision, by switching over to a Taylor series for small \delta and exploiting the expm1 function for larger \delta.
function g(δ::Union{Float64,ComplexF64})
if abs(δ) < 0.2
return 1 + δ*@evalpoly(δ^2, 1/3,-1/90,1/2520,-1/75600,1/2395008,-691/54486432000)
else
y = expm1(-δ)
return -2(y + δ*exp(-δ)) / y^2
end
end
(It’s possible that there is some clever algebraic transformation that lets you avoid the explicit Taylor series, but I’m not seeing it at the moment.)
16 Likes
you can also take the fact of the existence of expm1 and expanding each term:
f1(x)= 2*(1-x*exp(-x)-exp(-x))/(1-exp(-x))^2 #for comparison
function f2(x)
denom= -expm1(-x)
return 2/denom+ 2*x/denom- 2*x/(denom*denom)
end
Testing for equality on other number
julia> f1(4)
1.885271061051406
julia> f2(4)
1.8852710610514052
near zero:
julia> f1(1e-7)
0.9992008230547272
julia> f2(1e-7)
1.0000000335276127
julia> f2(1e-8)
1.0
julia> f1(1e-8)
0.0
to be fair it also fails, but around 1e-11
2 Likes
Often times for numerical computing, in the field I was in, we would either rigorously determine, or crudely determine limits for significance/guard digits. Rules of thumb for floats were usually 1e-5, for doubles usually 1e-10 (because we used ffts and things a lot).
Main take away is, there’s error in every calculation you do basically. There’s also a finite resolution to what a digital object can store/handle. When you do lots of calculations that error can propagate, or do unexpected things.
2 Likes
Does this means that in any serious calculations where human lives are at stake, you need to do the calculations twice? Once in Float64 and once in Float32. Then compare the results to make sure they come up with roughly the same value in case there are any numerical instabilities?
1 Like
Well I’m not so sure about the answer to that. Another thing that can be done, instead, would be to plot a range of solutions near the answer, and test results of known similar conditions.
I mostly did work like this. When something goes awry you’ll be able to tell most of the time. Complex real and imag values become underdetermined you’ll see steppy swaps. Really though you can either measure or find the floating point error for most calculations. Julia is actually amazing for this… I need to write a paper about it - but have negative time to do so. I have tons of results and code for it though…
When someone has done a careful analysis, this is not necessary—in some cases one can prove that the result has a certain accuracy.
3 Likes
thanks all of you for such detailed explanation and diverse solutions
Although it is not entirely foolproof, this method is likely to work well in practice. For a rather extensive discussion about such techniques, see for example this paper by Kahan:
The discussion about possible round-off error analysis schemes and Kahan’s opinion about them starts at page 13.
The first of the five schemes listed in this paper is the one you mention: running the same computation several times with different FP precisions. The second scheme (Kahan’s favorite) is very similar: it consists in running the same computation several times with the same FP precision, but with different rounding modes.
4 Likes
Sometimes numerical work may inform decisions and influence actions relevant to another’s health, wealth, or contentment. I have found it better to compute more accurately than may obtain using Float32 or Float64, then converting the likely more accurate result to Float32 or Float64. The choice of computational type and the choice of presentational type are made with an understanding of the specifics of the work and its context.
There are some publicly available Julia packages designed for this work. For your purposes, I suggest becoming familiar with DoubleFloats.jl. Where algorithmic choices, data uncertainties, or natural inexactness suggest characterizing the results as spans, The types ArbReal and ArbComplex from ArbNumerics.jl do just that. Also see packages from JuliaIntervals.
4 Likes
I saw your talk on this!!!
1 Like
You can also use TaylorSeries.jl to do Taylor series manipulations:
julia> using TaylorSeries
julia> n = 16
16
julia> t = Taylor1(Rational{Int}, n)
1//1 t + 𝒪(t¹⁷)
julia> ee = sum( (-t)^i / factorial(i) for i in 0:n )
1//1 - 1//1 t + 1//2 t² - 1//6 t³ + 1//24 t⁴ - 1//120 t⁵ + 1//720 t⁶ - 1//5040 t⁷ + 1//40320 t⁸ - 1//362880 t⁹ + 1//3628800 t¹⁰ - 1//39916800 t¹¹ + 1//479001600 t¹² - 1//6227020800 t¹³ + 1//87178291200 t¹⁴ - 1//1307674368000 t¹⁵ + 1//20922789888000 t¹⁶ + 𝒪(t¹⁷)
julia> 2 * (1 - t * ee - ee) / ((1 - ee)^2)
1//1 + 1//3 t - 1//90 t³ + 1//2520 t⁵ - 1//75600 t⁷ + 1//2395008 t⁹ - 691//54486432000 t¹¹ + 1//2668723200 t¹³ - 59513//156920924160000 t¹⁵ - 42397//94152554496000 t¹⁶ + 𝒪(t¹⁷)
(Note that although exp(t) is defined for a Taylor series, it currently converts the arguments to floats.)
You can also bound the remainder using TaylorModels.jl.
1 Like
That’s a neat writeup of very relevant problems and methods from a giant in the field, but it is surprising how jaded and bitter the tone of the article is.
1 Like | 2022-06-29T23:03:22 | {
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https://forum.math.toronto.edu/index.php?PHPSESSID=fd78vpp36m2806qcf102ad2st2&topic=2264.0;nowap | ### Author Topic: differences between proper and improper nodes (Read 7021 times)
#### Richard Qiu
• Newbie
• Posts: 1
• Karma: 0
##### differences between proper and improper nodes
« on: November 18, 2019, 02:19:33 PM »
Hello guys, could anyone help me to explain the differences between proper and improper nodes? btw, any suggestions on how to remember the types and stability of the critical points?
#### Amanda-fazi
• Newbie
• Posts: 3
• Karma: 1
##### Re: differences between proper and improper nodes
« Reply #1 on: November 18, 2019, 02:35:47 PM »
Since both proper and improper nodes have equal eigenvalues, the differences between these two nodes is that: proper node/star point has two independent eigenvectors, while improper/degenerate node has only one independent eigenvector by (A-rI)x =0, and we create a generalized eigenvector associated with the repeated eigenvalues by letting (A-rI)y = x.
« Last Edit: November 18, 2019, 03:21:26 PM by Amanda-fazi »
#### Amanda-fazi
• Newbie
• Posts: 3
• Karma: 1
##### Re: differences between proper and improper nodes
« Reply #2 on: November 18, 2019, 02:49:15 PM »
There are mainly 5 cases of Eigenvalues(from book Elementary Differential Equations and Boundary Value Problems-11th Edition section 9.1):
as it is mentioned above, the equal eigenvalues case mentioned above is CASE 3.
CASE 1: Real, Unequal Eigenvalues of the Same Sign
CASE 2: Real Eigenvalues of Opposite Sign ->saddle point
CASE 3: Equal Eigenvalues
CASE 4: Complex Eigenvalues with Nonzero Real Part
CASE 5: Pure Imaginary Eigenvalues ->center
After memorized there are five cases, CASE 1, CASE 3 and CASE 4 have two branches while the rest of the cases(CASE 2 and CASE 5) only have one:
to be more specific:
CASE 1: Real, Unequal Eigenvalues of the Same Sign separated into:
a)lambda1 >lambda2 >0:
critical point called node/nodal source
a)lambda1 <lambda2 <0:
critical point called node/nodal sink
CASE 3:Equal Eigenvalues separated into:
a)two independent eigenvectors:
critical point called proper node or star point
b)one independent eigenvector:
critical point called improper node or degenerate node
CASE 4:Complex Eigenvalues with Nonzero Real Part separated into:
a)pointing-outward trajectories as lambda > 0:
critical point called spiral source
a)pointing-inward trajectories as lambda < 0:
critical point called spiral sink
For the stability, as long as there is one lambda>0, then it is unstable, and the last one lambda=0 is stable. For the rest of them, asymptotically stable applied.
« Last Edit: November 18, 2019, 03:20:51 PM by Amanda-fazi »
#### anntara khan
• Jr. Member
• Posts: 6
• Karma: 0
##### Re: differences between proper and improper nodes
« Reply #3 on: November 20, 2019, 04:03:25 PM »
I made this handy color coded guide to help me remember all the cases:
« Last Edit: December 05, 2019, 09:37:34 PM by anntara khan »
#### anntara khan
• Jr. Member
• Posts: 6
• Karma: 0
##### Re: differences between proper and improper nodes
« Reply #4 on: December 11, 2019, 09:59:34 PM »
Based on the stability near locally linear system I have extended the previously posted table, hope this helps remembering | 2022-12-04T19:31:29 | {
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https://www.thelaunchpadtech.com/r52a72/g6p3jm.php?c9f433=conjugate-transpose-of-a-matrix-example | Here $*$ denotes the conjugate transpose. Matrix representation. In mathematics, the conjugate transpose or Hermitian transpose of an m-by-n matrix $\boldsymbol{A}$ with complex entries is the n-by-m matrix $\boldsymbol{A}^\mathrm{H}$ obtained from $\boldsymbol{A}$ by taking the transpose and then taking the complex conjugate of each entry. Returns the (complex) conjugate transpose of self.. real part of the matrix component and the second element of each pair is the imaginary part of the corresponding matrix component. With the help of Numpy numpy.matrix.getH() method, we can make a conjugate Transpose of any complex matrix either having dimension one or more than more.. Syntax : matrix.getH() Return : Return conjugate transpose of complex matrix Example #1 : In this example we can see that with the help of matrix.getH() we can get the conjugate transpose of a complex matrix having any dimension. Example.' Are there other cases when a matrix commutes with its transpose ? For example, if B = A' and A(1,2) is 1+1i, then the element B(2,1) is 1-1i. The complex conjugate transpose of a matrix interchanges the row and column index for each element, reflecting the elements across the main diagonal. does not affect the sign of the imaginary parts. The conjugate transpose of a matrix with real entries reduces to the transpose of , as the conjugate of a real number is the number itself. Thanks for contributing an answer to Mathematics Stack Exchange! The conjugate transpose of a matrix can be denoted by any of these symbols: â, commonly used in linear algebra numpy.matrix.H¶ matrix.H¶. example. In all common spaces, the conjugate and transpose operations commute i.e., A H ⦠The operation also negates the imaginary part of any complex numbers. Definition. Transpose of a linear mapping. The complex conjugate transpose of a matrix interchanges the row and column index for each element, reflecting the elements across the main diagonal. Hence, instead of storing all entries of the matrix, Sage only stores the interesting entries in a dictionary, hence the name sparse matrix. The matrix in Example 23 is invertible, and the inverse of the transpose is the transpose of the inverse. It is easy to verify cX*cX' = sum(abs(cX)^2), where cX' is the conjugate transpose. (The complex conjugate of ⦠In mathematics, the conjugate transpose, Hermitian transpose, Hermitian conjugate, or adjoint matrix of an m-by-n matrix A with complex entries is the n-by-m matrix A * obtained from A by taking the transpose and then taking the complex conjugate of each entry (i.e., negating their imaginary parts but not their real parts). Annihilator. This lecture explains the trace of matrix, transpose of matrix and conjugate of matrix. To understand the properties of transpose matrix, we will take two matrices A and B which have equal order. For example, if A(3,2) is 1+2i and B = A. In mathematics, the conjugate transpose or Hermitian transpose of an m-by-n matrix A with complex entries is the n-by-m matrix A * obtained from A by taking the transpose and then taking the complex conjugate of each entry (i.e., negating their imaginary parts but not their real parts). $\begingroup$ I got the conjugate. Let V be an abstract vector space over a field F. A functional T is a function T:V â F that assigns a number from field F to each vector x ε V. Def. Hermitian conjugate) of a vector or matrix in MATLAB. does not affect the signal of the imaginary parts. ', then the component B(2,3) is also 1+2i. If the conjugate transpose of a square matrix is equal to its inverse, then it is a unitary matrix. returns the nonconjugate transpose of A, that is, interchanges the row and column index for each element.If A contains complex elements, then A.' # Good! I'm not sure at all how to convert the complex conjugate transform to c, I just don't understand what that line does. Dual space, conjugate space, adjoint space. Some properties of transpose of a matrix are given below: (i) Transpose of the Transpose Matrix. Examples. For a square matrix A it is the matrix . Conjugate Transpose of Real Matrix; The complex conjugate transpose of a matrix interchanges the row and column ctranspose and transpose produce the, Operations with Matrices ! In mathematics, the conjugate transpose or Hermitian transpose of an m-by-n matrix A with complex entries is the n-by-m matrix Aâ obtained from A by taking the transpose and then taking the complex conjugate of each entry. In all common spaces (i.e., separable Hilbert spaces), the con Other names for the conjugate transpose of a matrix are Hermitian conjugate, bedaggered matrix, adjoint matrix or transjugate. (The complex conjugate of a + bi, where a and b are reals, is a â bi.) I have to further multiply 1x4 matrix with 4x1 matrix and get a scalar. To find the conjugate trans-pose of a matrix, we first calculate the complex conjugate of each entry and then take the transpose of the matrix, as shown in the following example. Calculates the conjugate matrix. Remember that the complex conjugate of a matrix is obtained by taking the complex conjugate of each of its entries (see the lecture on complex matrices). B = A.' This is equivalent to Conj(t.default(x)). ', there is a period in front of the apostrophe. Theorems. Conjugate Transpose for Complex Matrix. The transpose of the conjugate of a matrix. Motivation The conjugate transpose can be motivated by noting that complex numbers can be usefully represented by 2×2 real matrices, obeying matrix ⦠| 2021-09-27T22:15:19 | {
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https://math.stackexchange.com/questions/4172776/probability-of-rolling-3-dice-numbers-which-product-is-a-multiple-of-6/4172797 | probability of rolling 3 dice numbers which product is a multiple of 6
As the question states, I need to find the probability of rolling 3 dice numbers which product is a multiple of 6. what I have tried:
$$\omega={(6,k,k), (3,2,k), (3,4,k))}$$
$$P(A) = \frac{3\times(1^1\times6^2) + 6\times(1^1\times1^1\times6^1) + 6\times(1^1\times1^1\times6^1)}{6^3}=\frac{180}{216}=\frac{5}{6}$$
I know writing $$1^1$$ seems silly, but just trying to show what I'm doing
Thanks for helping and sorry if it is a bad question
• This looks wrong at a glance as it looks like you are overcounting. For instance, where does the outcome where all three dice show a $6$ get counted and how many times did you count it? It should only have been counted once, but it looks like you counted it three times. What about the outcome where you rolled a $2,3,6$? in some order? Was it counted in the first category? or in the second? Jun 14 at 14:41
• You have overcounted since there is an overlap between the cases you describe. E.g. (6,2,3) is in two of the cases so is counted twice. Jun 14 at 14:42
• I think it would be easier to count the outcomes where the product is not even, count the outcomes where the product is not a multiple of three, and count the outcomes where the product is neither even nor a multiple of three, and use that to proceed to a solution. Jun 14 at 14:43
• @JMoravitz what if i use the same way i did but substract the cases where i overcounted? Jun 14 at 14:47
• That seems tedious. The question is crafted in such a way as to specifically suggest using inclusion-exclusion and basic facts about divisibility to approach. If you want to explore other additional approaches after the fact, go ahead and revisit the problem, but I strongly encourage you to use the approach I suggested in my second comment first. Jun 14 at 14:48
Let the dice be thrown in sequence or otherwise be distinguishable some other way. There are $$6^3 = 216$$ total different equally likely outcomes.
Approach by trying to count how many of these outcomes were "bad" and should be removed because they either are not multiples of $$2$$, are not multiples of $$3$$, and keeping in mind that they might have been both.
Of these, $$3^3=27$$ are "bad" because their product of the results is not even (and thus not a multiple of six), seen by noting a product is odd iff all terms in the product is odd, there being three odd numbers.
$$~$$
Similarly, $$4^3=64$$ are "bad" because their product of results is not a multiple of three (and thus not a multiple of six), seen by noting a product is not a multiple of three iff all terms in the product are not multiples of three, there being four possible non-multiple of three numbers to choose from here.
$$~$$
Finally, some of these "bad" outcomes we counted twice as they were simultaneously not even and not multiples of three. These would be the outcomes consisting only of $$1$$'s and $$5$$'s, there being $$2^3=8$$ of which. We keep this in mind to correct our count as we proceed with inclusion-exclusion.
The probability then is:
$$\frac{6^3-3^3-4^3+2^3}{6^3} = \frac{216 - 27 - 64 + 8}{216} = \frac{133}{216}$$
Let $$a_2,b_2,c_2$$ ($$a_3,b_3,c_3$$) denote the exponents of $$2$$ ($$3$$) in the prime factorization of the number shown on the three dice. Then the probability you want is
$$P(a_2+b_2+c_2>0, a_3+b_3+c_3>0) = 1-P(a_2+b_2+c_2=0\ or \ a_3+b_3+c_3=0)$$
$$= 1-P(a_2+b_2+c_2=0)-P(a_3+b_3+c_3=0)+P(a_2+b_2+c_2=0, a_3+b_3+c_3=0)$$
$$(by\ independence)= 1-P(a_2=0)^3-P(a_3=0)^3+P(a_2=0,a_3=0)^3$$
$$= 1-P(1,3,\ or \ 5)^3-P(1,2,4, \ or \ 5)^3+P(1 \ or \ 5)^3$$
$$= 1 - (\frac{3}{6})^3 - (\frac{4}{6})^3+(\frac{2}{6})^3 = \frac{133}{216}.$$
• very elegant approach ! Jun 14 at 15:32 | 2021-09-20T15:22:03 | {
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https://math.stackexchange.com/questions/1265071/which-matrices-are-covariances-matrices | # Which matrices are covariances matrices?
Let $V$ be a matrix.
What conditions should we require so that we can find a random vector $X = (X_1, \dots, X_n)$ so that $V = Var(X)$?
Of course necessary conditions are:
• All the elements on the diagonal should be positive
• The matrix has to be symmetric
• $v_{ij} \le \sqrt{v_{ii}v_{jj}}$ (Because of $Cov(X_i, X_j) \le \sqrt{Var(X_i) Var(X_j)})$
But I am sure these are not sufficient as I have a counterexample.
So what other properties we should require on a matrix so that it can be considered a covariance matrix?
I think I cleared this up sufficiently.
Okay, so
1) If $V$ is not semi definite positive, then such a vector $X$ does not exists. (Since all covariances matrix are semi definite positive)
2) If $V$ is symmetric semidefinite positive, then such an $X$ exists! [0]
This implies that
$$\text{exists a random vector X: V = Cov(X)} \iff \text{V is symmetric positive semidefinite}$$
Since we know that those I listed in the question are necessary condition for $V$, we deduce that all symmetric semidefinite positive matrices have elements on the diagonal $\ge 0$ and are such that $v_{ij} \le \sqrt{v_{ii}v_{jj}}$.
These are not sufficient though for a matrix to be semidefinite positive but sufficient conditions are well known, after all.
[0] Proof
Since $V$ is symetric is possible to find an orthogonal matrix $Q$ such that $V = QDQ^T$, where $D$ is a diagonal matrix whose values are the eigenvalues of $V$. If $V$ is semipositive definite the elements of $D$ are all $\ge 0$, hence we can find $X$ such that $D = Cov(X)$ (just take all the variables independent with the specified variance)
It follows that the random vector $QX$ has covariance equal to $$Cov(QX) = QCov(X)Q^T = QDQ^T = V$$
• There's something unclear about your proof: where exactly did you use the condition $v_{ij} \leq \sqrt{v_{ii}v_{jj}}$? Is this a consequence of the positive semidefinite/symmetric matrix condition? Beacuse if not, you there's something wrong in your conclusion. – RandomGuy Oct 31 '16 at 14:32
• @RandomGuy I didn't. I showed that if a matrix is symmetric, positive semidefinite, then it can be written as $Cov(X)$; therefore it must hold that $v_{ij} \le \sqrt{v_{ii}v_{j}}$ for every symmetric semidefinite positive matrix, because it holds for $Cov(X)$. You can consider this a proof of this statement. Of course, if the statement is false, then my proof is invalid :) – Ant Nov 1 '16 at 20:04
• yes, now it makes sense. My question was about the fact that the statement of your original question included that inequality as a condition and in the answer you answer your own question, which seemed to imply that you were using that condition as a hypothesis. – RandomGuy Nov 2 '16 at 8:18 | 2019-07-22T03:22:42 | {
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https://math.stackexchange.com/questions/2403597/area-of-overlapping-quadrilateral-tiles/2403618 | # Area of Overlapping Quadrilateral Tiles
I was studying for some quizzes when I stumbled across this question. It goes like this:
Two regular quadrilaterals vinyl tiles each of $1$ foot long on each sides overlap each other such that the overlapping region is a regular octagon. What is the area of the overlapping region?
My work:
I imagined the problem like this:
The area of the shaded blue region is what I want to get. I don't know if one side of the octagon is one-third of that of the side of square, because my calculated area is slightly larger than that of my book's answer.
How do you get the area of the blue region above?
Let $a$ be a length-side of the octagon $ABCDE...$, where $BC$ placed on the square side.
Thus, since $\sin45^{\circ}=\frac{1}{\sqrt2}$, we obtain:$$\frac{1}{2}-\frac{a}{2}=a\cdot\frac{1}{\sqrt2},$$ which gives $a=\sqrt2-1$.
Thus, $\frac{1}{2}-\frac{a}{2}=\frac{1}{2}-\frac{\sqrt2-1}{2}=\frac{2-\sqrt2}{2}$ and the needed area is $$1-4\cdot\frac{1}{2}\left(\frac{2-\sqrt2}{2}\right)^2=1-\frac{1}{2}(2-\sqrt2)^2=2\sqrt2-2$$
• How did you derived the expression above? What is the thought process? – Palautot Ka Aug 23 '17 at 15:16
• @Palautot Ka Say me please, are my answer and the answer in your book, same? – Michael Rozenberg Aug 23 '17 at 15:35
• The area of a regular polygon is $A = \frac{a^2 n }{4} \cot \frac{180^o}{a}.$ Using your $a = \sqrt{2} -1$ and $n = 8$ sides, I got the area to be 0.8284 square feet. Is your value $4 \sqrt{2} - 5$ the area of the blue octagon? – Palautot Ka Aug 23 '17 at 15:35
• @Palautot Ka I fixed my post. See now please. – Michael Rozenberg Aug 23 '17 at 15:39
• The book's answer is 119.29 square inchesXD – Palautot Ka Aug 23 '17 at 15:40
The area of the shaded region will be an octagon - one square minus 4 white triangles. Say the legs of the right triangles are l. Then $$2l+{\sqrt{2l^2}} = 1$$ from the Pythagorean theorem. Solving for l we get $$1-{\sqrt2}/2$$ The area will then be $$1-4*l^2/2$$ giving us $$2*{\sqrt2}-2$$
• I like the insight...... – Palautot Ka Aug 23 '17 at 15:45
The height of a triangle is the half of the diagonal of a square minus the side of a square.
The area of a triangle (isoceles rectangle) is the square of the height.
Hence the blue area
$$1^2-4\frac{\left(\sqrt2-1\right)^2}4=2\sqrt2-2.$$ | 2019-12-15T08:17:12 | {
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https://math.stackexchange.com/questions/1207829/how-to-prove-that-sqrt3-2-sqrt3-4-is-irrational | How to prove that $\sqrt[3] 2 + \sqrt[3] 4$ is irrational? [duplicate]
So while doing all sorts of proving and disproving statements regarding irrational numbers, I ran into this one and it quite stumped me:
Prove that $\sqrt[3]{2} + \sqrt[3]{4}$ is irrational.
I tried all the usual suspects like playing with $\sqrt[3]{2} + \sqrt[3]{4} = \frac{a}{b}$ for $a,b\in \mathbb{Z}$ , but got nowhere.
I also figured maybe I should play with it this way:
$2^\frac{1}{3} + 4^\frac{1}{3}=2^\frac{1}{3} + (2^2)^\frac{1}{3}=2^\frac{1}{3} + 2^\frac{2}{3}=2^\frac{1}{3} + 2^\frac{1}{3}\times 2^\frac{1}{3}=2^\frac{1}{3}(1+2^\frac{1}{3})$
But there I got stumped again, because while $1+2^\frac{1}{3}$ is irrational, nothing promises me that $2^\frac{1}{3} \times (1+2^\frac{1}{3})$ is irrational, and I feel like trying to go further down this road is moot.
So what am I missing (other than sleep and food)? What route should I take to prove this? Thanks in advance!
marked as duplicate by Watson, Chinnapparaj R, user10354138, Brahadeesh, KReiserNov 26 '18 at 8:38
• – punctured dusk Apr 7 '15 at 12:22
• – Watson Nov 25 '18 at 19:45
• – Watson Nov 25 '18 at 21:45
• – Watson Jan 29 at 9:06
Note
$$1 + \sqrt[3]{2} + \sqrt[3]{4} = \frac{(\sqrt[3]{2})^3 - 1}{\sqrt[3]{2} - 1} = \frac{1}{\sqrt[3]{2} - 1}.$$
So if $\sqrt[3]{2} + \sqrt[3]{4}$ is rational, then $1/(\sqrt[3]{2} - 1)$ is rational, which implies $\sqrt[3]{2} - 1$ is rational. Then $\sqrt[3]{2}$ is rational, a contradiction.
• Very nice answer. +1 – Timbuc Mar 26 '15 at 18:00
• Thanks for the reply. I must be tired or confused because I'm missing something very basic and don't realize how $1 + \sqrt[3]{2} + \sqrt[3]{4} = \frac{(\sqrt[3]{2})^3 - 1}{\sqrt[3]{2} - 1}$. I apologize for my stupidity - it's been a tough day. – Elad Avron Mar 26 '15 at 18:01
• Let $x = \sqrt[3]{2}$. Use the factorization $(1 + x + x^2)(x - 1) = x^3 - 1$ to obtain the identity. – kobe Mar 26 '15 at 18:04
• Of course, silly me. Thank you so much, this solution is brilliant! – Elad Avron Mar 26 '15 at 18:09
• @Elad One easly proves a much more general result: for irrational cube roots of rationals, said property fails only for $\,\sqrt[3]1,\,$ see my answer. $\ \$ – Bill Dubuque Mar 26 '15 at 20:21
$y = \sqrt[3]{2} + \sqrt[3]{4}$ is a root of the equation $y^3 - 6y - 6 = 0$.
(To see this, let $x = \sqrt[3]{2}$ and $y = \sqrt[3]{2}+\sqrt[3]{4}=x+x^2$. Then $y^3 = x^3 + 3x^4 + 3x^5 + x^6 = 2 + 6x + 6x^2 + 4 = 6 + 6y$.)
By the rational root theorem, we know that any rational roots of $y^3 - 6y - 6 = 0$ would have to be in the set $\{\pm1,\pm2,\pm3,\pm6\}$; we can quickly confirm that none of these is in fact a root, meaning that the equation does not have any rational roots.
Therefore, $y$ must be irrational.
• Short, sweet, and intuitive. Very nice. – Jon Mar 26 '15 at 21:16
An easy approach: If $p=\sqrt[3]{2}+\sqrt[3]{4}$ is rational:
$$p^2 = \sqrt[3]{4}+2\cdot 2+ 2\sqrt[3]{2} = p+4+\sqrt[3]{2}.$$
So $\sqrt[3]{2}=p^2-p-4$ would be rational.
An alternative, more general approach.
Claim: If $a,b$ are integers that are not perfect cubes, and $a\neq -b$, then $\sqrt[3]{a}+\sqrt[3]b$ is irrational.
Proof:
Assume $\sqrt[3]{a}+\sqrt[3]{b}$ is rational. Cube it and get:
$$(\sqrt[3]{a}+\sqrt[3]{b})^3 = a+ 3\sqrt[3]{ab}(\sqrt[3]{a}+\sqrt[3]{b}) + b$$
Now, since $a,b$ are rational and $\sqrt[3]{a}+\sqrt[3]{b}$ is non-zero and rational, this means that $\sqrt[3]{ab}$ is rational.
Letting $p=\sqrt[3]{a}+\sqrt[3]{b}$ and $q=\sqrt[3]{ab}$, this means that
$$(x-\sqrt[3]{a})(x-\sqrt[3]{b}) = x^2-px+q$$ is a rational polynomial. It shares at least one root with $x^3-b$, But the GCD of these two polynomials has to be a rational polynomial, so the GCD cannot be linear (since it would be $x-\sqrt[3]b$, which is not a rational polynomial.)
This means that $x^2-px+q$ has to divide $x^3-b$. That means that $\sqrt[3]{a}$ is a root of $x^3-b$, which means that $a=b$. But there are no repeated roots of $x^3-b$, which yields a contradiction.
Corollary: If $a,b$ are rationals such that $a\neq -b$ and $\sqrt[3]{a}$ and $\sqrt[3]{b}$ are irrational, then $\sqrt[3]{a}+\sqrt[3]{b}$ is irrational.
Proof: Rationalize the denominators and revert to the above theorem for integers.
It holds true for any irrational $\,x=\sqrt[3]n,\,$ except $\,n=1\,$ (so $\,x^2+x=-1\in\Bbb Q)$
More generally: if $\ r\in\Bbb Q\$ and $\,x=\sqrt[3]r\not\in\Bbb Q\,$ then $\,x^2+x = q\in\Bbb Q\iff r = 1.$
Proof $\,\ qx = x^3+x^2 = r+x^2\$ so $\ qx-r = x^2 = q-x,\$ so $\,(\color{#c00}{q\!+\!1})\,x = r+q.$
Therefore $\,\ x\not\in\Bbb Q\,\Rightarrow\,\color{#c00}{q = -1}\,\Rightarrow\, 0 = (\color{#c00}{q\!+\!1})x = r+q = r-1,\$ thus $\,\ r = 1.$
• Um, if $n=1$, isn't it $x^2+x=2$? And is $x=\sqrt[3]{1}$ irrational? I'm confused. – Thomas Andrews Mar 29 '15 at 4:54
• @Thomas The hypothesis is that $\,x\,$ is an irrational cube root of $\,n.\ \$ – Bill Dubuque Mar 29 '15 at 6:41
• Well, $\sqrt[3]{n}$ is a single-valued function, by convention. You could say any irrational $x$ with $x^3=n$, I suppose. – Thomas Andrews Mar 29 '15 at 15:25
• @Thomas Algebraists not too infrequently overload the notation to denote any root. I will change it to avoid possible confusion. Thanks for pointing that out. – Bill Dubuque Mar 29 '15 at 16:25
Here is another approach which generalises to many similar examples, and which involves no complicated simplifications of surds. (In fact, no easy simplifications of surds either.)
First, $\sqrt[3]2$ is an algebraic integer, that is, it is a root of $$x^3-2$$ which is a monic polynomial (leading coefficient $1$) with integer coefficients. Similarly, $\sqrt[3]4$ is an algebraic integer.
It is known that the sum of two algebraic integers is an algebraic integer. Thus, $x=\sqrt[3]2+\sqrt[3]4$ is an algebraic integer.
It is also known that if an algebraic integer is rational, then it is a rational integer - that is, an ordinary integer, $0,1,-1,2,-2$ etc. However, it is not hard to find the estimates $$1<\sqrt[3]2<\frac43 ,\quad \frac43<\sqrt[3]4<\frac53\ ;$$ so $\frac73<x<3$, hence $x$ is not an integer and must be irrational.
If $a$ is rational, $a^2-a-4$ is rational. What is $a^2-a-4$ when $a=\sqrt[3]2+\sqrt[3]4$?
(After you work that out, I bet you'll wonder where I got $a^2-a-4$ from. Or, try to figure it out yourself. Now that you have a sort of idea on how to deal with this sort of problem, I leave you with another, similar problem: What polynomial can I use to prove that $\sqrt[3]3+\sqrt[3]9$ is irrational?)
• I was going to suggest proving $\sqrt[3]2+\sqrt[3]3$ irrational, but then I realized that the polynomial involved is an eight-degree one. And I'm not sure how hard the one I gave is; I didn't try it yet. – Akiva Weinberger Mar 26 '15 at 20:49 | 2019-06-17T13:25:31 | {
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https://math.stackexchange.com/questions/3307711/whats-my-mistake-in-finding-int-0-infty-dx-e-ax2-sinb-x2 | # What's my mistake in finding $\int_0^\infty dx e^{-ax^2} \sin(b/x^2)$?
I want to evaluate $$I=\int_0^\infty dx e^{-ax^2} \sin(b/x^2)$$ for $$a,b>0$$. A first simplification is to substitute $$y=x/\sqrt{a}$$ and define $$c=ab>0$$ to obtain $$I=\frac{1}{\sqrt{a}} \int_0^\infty e^{-x^2} \sin(c/x^2)$$ Now my idea was to use the Taylor series for the sine $$I=a^{-1/2} \int_0^\infty dx e^{-x^2} \sum_{k=0}^\infty \frac{(-1)^k}{(2k+1)!}\left(\frac{c}{x^2} \right)^{2k+1}$$ Now I interchange sum and integral although I have no justification $$I=a^{-1/2} \sum_{k=0}^\infty \frac{(-1)^kc^{2k+1}}{(2k+1)!} \int_0^\infty dx e^{-x^2} \left(\frac{1}{x^2} \right)^{2k+1}$$ Substituting $$t=x^2$$ in the integral we obtain the gamma function $$I=a^{-1/2} \sum_{k=0}^\infty \frac{(-1)^kc^{2k+1}}{(2k+1)!} \frac{1}{2} \Gamma(-2k-1/2)$$ Using $$\Gamma \left({\frac{1}{2}}-n\right)={(-4)^{n}n! \over (2n)!}{\sqrt {\pi }}$$ (which can be shown using the reflection formula and the duplication formula for the Gamma function) with $$n=2k+1$$ I obtain $$I=a^{-1/2} \sqrt{\pi} \sum_{k=0}^\infty \frac{(-1)^k(-4c)^{2k+1}}{(4k+2)!} \frac{1}{2}$$ or $$I=-\frac{1}{2}\sqrt{\frac{\pi}{a}} \sum_{k=0}^\infty \frac{(-1)^k(4c)^{2k+1}}{(4k+2)!}=-\frac{1}{2}\sqrt{\frac{\pi}{a}} \sin(\sqrt{2c}) \sinh(\sqrt{2c})$$ where I used wolfram alpha for the last series.
The problem: The above result is wrong. It should be (wolfram alpha and Gradshteyn) $$I=\frac{1}{2}\sqrt{\frac{\pi}{a}} \sin(\sqrt{2c})\exp(-\sqrt{2c})$$
The question: Can someone spot my mistake? Was it interchanging the limits? I would also be interested in your solutions to the integral $$I$$ using other approaches.
• After the interchange, all the integrals diverge at $0$, aren't they? – user58697 Jul 29 at 19:58
• @user58697 Yes, you are right. Probably that is the mistake... Do you have an idea for an alternative approach to the original integral? – thomasfermi Jul 29 at 20:09
The Glasser's master theorem is a useful tool for the solution. First, use Euler's formula to decompose the sine term into the sum of exponentials. Then it boils down to computing the integral of the form
$$J(p) = \int_{0}^{\infty} \exp\left( -a x^2 - \frac{p}{x^2} \right) \, \mathrm{d}x.$$
Assume for a moment that $$a, p > 0$$. Then by completing the square, we get
$$J(p) = \int_{0}^{\infty} \exp\left( -a \left( x - \frac{\smash{\sqrt{p/a}}}{x} \right)^2 - 2\sqrt{ap} \right) \, \mathrm{d}x.$$
Then by the Glasser's master theorem and the gaussian integral, this evaluates to
$$J(p) = \int_{0}^{\infty} \exp\left( -a x^2 - 2\sqrt{ap} \right) \, \mathrm{d}x = \frac{1}{2}\sqrt{\frac{\pi}{a}} \exp(-2\sqrt{ap}). \tag{*}$$
Although $$\text{(*)}$$ is originally proved for $$p > 0$$, both sides of $$\text{(*)}$$ define holomorphic functions for $$p$$ in the right-half plane $$\mathbb{H}_{\to} = \{z \in \mathbb{C} : \operatorname{Re}(z) > 0\}$$ and are continuous on the closed right-half plane $$\overline{\mathbb{H}_{\to}}$$. So by the identity theorem and continuity, $$\text{(*)}$$ extends to all of $$p \in \overline{\mathbb{H}_{\to}}$$. In particular, plugging $$p = \pm ib$$ for $$b > 0$$, we get
$$J(\pm ib) = \frac{1}{2}\sqrt{\frac{\pi}{a}} \exp(-2\sqrt{\pm i ab}) = \frac{1}{2}\sqrt{\frac{\pi}{a}} \exp(-\sqrt{2c}(1\pm i)).$$
Therefore
$$I = \frac{J(-ib) - J(ib)}{2i} = \frac{1}{2}\sqrt{\frac{\pi}{a}} \sin(\sqrt{2c}) \exp(-\sqrt{2c}).$$
• Thanks for your very nice solution. I had never heard of Glasser's theorem before! – thomasfermi Jul 30 at 22:03 | 2019-09-19T18:55:44 | {
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http://mathhelpforum.com/algebra/86584-help-finding-sum-series-print.html | # help finding sum of series
Printable View
• April 30th 2009, 01:11 AM
adhyeta
help finding sum of series
got 2 series.
1.
12 - 22 + 32 - 42 + ................. 992 - 1002
(ok those are just squares)
2.
1.1! + 2.2! + ....... 50.50!
need to find sum...
• April 30th 2009, 03:07 AM
pickslides
Quote:
Originally Posted by adhyeta
got 2 series.
1.
12 - 22 + 32 - 42 + ................. 992 - 1002
(ok those are just squares)
need to find sum...
$1^2-2^2+3^2-4^2+\cdots+99^2-100^2$
$1-4+9-16+25-36+\cdots+99^2-100^2$
finding the difference in pairs now giving 50 terms.
$-3-7-11-\cdots$
Use sum of an arithmetic sequence
$S_n = \frac{n}{2}(2a+(n-1)d)$
where n = number of terms = 50
a = the first term = -3 and d = the common difference between terms = -4
$S_{50} = \frac{50}{2}(2(-3)+(50-1)(-4)) = \cdots$
• April 30th 2009, 04:22 AM
Soroban
Hello, adhyeta!
Are you familiar with these summation formulas?
. . $\sum^n_{k=1} 1 \;\;= \;n$
. . $\sum^n_{k=1}k \;\;=\;\frac{n(n+1)}{2}$
Quote:
$1)\;\;S \;=\;1^2 - 2^2 + 3^2 - 4^2 + \hdots + 99^2 - 100^2$
$\text{We have: }\;S \;=\;\underbrace{(1^2 + 3^2 + 5^2 + \hdots + 99^2)}_{\text{first 50 odd squares}} \;-\; \underbrace{(2^2+4^2+6^2 + \hdots + 100^2)}_{\text{first 50 even squares}}$
. . $\text{Then: }S \;=\qquad\qquad \sum^{50}_{k=1}(2k-1)^2 \qquad - \qquad\qquad \sum^{50}_{k=1}(2k)^2$
Hence: . $S \;=\;\sum^{50}_{k=1}\bigg[(2k-1)^2 - (2k)^2\bigg] \;=\;\sum^{50}_{k=1}(1-4k)$
. . $= \;\;\sum^{50}_{k=1}\!1 \;-\; 4\!\sum^{50}_{k=1}k \;\;=\;\;50 - 4\,\frac{50\cdot51}{2} \;\;=\;\;\boxed{-5,\!050}$
• April 30th 2009, 06:10 AM
adhyeta
Quote:
Originally Posted by pickslides
$1^2-2^2+3^2-4^2+\cdots+99^2-100^2$
$1-4+9-16+25-36+\cdots+99^2-100^2$
finding the difference in pairs now giving 50 terms.
$-3-7-11-\cdots$
Use sum of an arithmetic sequence
$S_n = \frac{n}{2}(2a+(n-1)d)$
where n = number of terms = 50
a = the first term = -3 and d = the common difference between terms = -4
$S_{50} = \frac{50}{2}(2(-3)+(50-1)(-4)) = \cdots$
thats not coming out to be correct...its -5050.
refer to soroban's solution.
• April 30th 2009, 02:11 PM
pickslides
Quote:
Originally Posted by adhyeta
thats not coming out to be correct...its -5050.
refer to soroban's solution.
$
S_{50} = \frac{50}{2}(2(-3)+(50-1)(-4))$
$
S_{50} = 25(-6+(49)(-4))$
$
S_{50} = 25(-6+-196)$
$
S_{50} = 25(-202)$
$
S_{50} = -5050$
Worked for me!
• April 30th 2009, 03:04 PM
Jester
Quote:
Originally Posted by adhyeta
got 2 series.
1.
12 - 22 + 32 - 42 + ................. 992 - 1002
(ok those are just squares)
2.
1.1! + 2.2! + ....... 50.50!
need to find sum...
Here's the second one. Let
$S_n = 1 \cdot 1! + 2 \cdot 2! + 3 \cdot 3! + \cdots + (n-1) \cdot (n-1)! + n \cdot n!$
$= (2-1) 1! + (3-1)2! + (4-1)3! + \cdots + (n -1) \cdot (n-1)! + (n+1 -1) \cdot n!
$
$= (2! - 1!) + (3! - 2!) + (4! - 3!) + \cdots + (n! - (n-1)!) + ((n+1)! - n!)$
$
= (n+1)! - 1
$
so here $n = 50$ so the answer $S = 51!-1$.
• April 30th 2009, 05:16 PM
adhyeta
Quote:
Originally Posted by pickslides
$
S_{50} = \frac{50}{2}(2(-3)+(50-1)(-4))$
$
S_{50} = 25(-6+(49)(-4))$
$
S_{50} = 25(-6+-196)$
$
S_{50} = 25(-202)$
$
S_{50} = -5050$
Worked for me!
hey! so sorry. i think i went wrong with the calc.(-calculation-)
(Happy)
• May 3rd 2009, 07:16 AM
adhyeta
Quote:
Originally Posted by danny arrigo
Here's the second one. Let
$S_n = 1 \cdot 1! + 2 \cdot 2! + 3 \cdot 3! + \cdots + (n-1) \cdot (n-1)! + n \cdot n!$
$= (2-1) 1! + (3-1)2! + (4-1)3! + \cdots + (n -1) \cdot (n-1)! + (n+1 -1) \cdot n!
$
$= (2! - 1!) + (3! - 2!) + (4! - 3!) + \cdots + (n! - (n-1)!) + ((n+1)! - n!)$
$
= (n+1)! - 1
$
so here $n = 50$ so the answer $S = 51!-1$.
hey! can we similarily prove the sum $n^{2}.n!$??? | 2016-05-24T09:01:45 | {
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https://dsp.stackexchange.com/questions/60939/calculating-amplitude-time-shift-and-phase-of-third-sinusoid-as-a-sum-of-two-ot | Calculating amplitude, time-shift and phase of third sinusoid as a sum of two others
My apologies in advance, I am new to this digital signal processing. Question is at the bottom
I have two 4000hz sinusoids given by the following formulas: unknown frequency such that: $$x_1(t) = 24\cos(2\pi4000(t-t_{m1})) \\ x_2(t) = 28.8 \cos(2\pi4000(t-t_{m2}))$$ and a third given by $$x_3(t) = x_1(t) + x_2(t)$$ where $$t_{m1} = 0.000775 \\ t_{m2} = -0.00258$$
I have calculated the phase of the first two sinusoids (in radians) as $$\phi_1 = 0.329 \\ \phi_2 = -1.132$$ with the formula $$\phi_i = \omega_if_i$$ where $$\omega_i = 2\pi f$$
Finally, I need to calculate the amplitude $$A_3$$ and phase $$\phi_3$$ of $$x_3(t)$$. I am unsure which formula to use to do so. I believe $$A_3$$ is calculated as: $$a_3 = \sqrt{A_1^2 + A_2^2} = 37.49$$
How would I calculate the time shift and phase angle of $$x_3(t)$$?
• Are you familiar with complex numbers ? – Hilmar Sep 27 '19 at 18:38
• @hilmar I am, but not with respect to this topic. – kp-a Sep 27 '19 at 20:34
Here's how I would do it.
Use the Cosine angle addition formula on $$x_1$$ and $$x_2$$:
$$\cos( \alpha + \beta) = \cos( \alpha ) \cos(\beta)-\sin( \alpha ) \sin(\beta)$$
Like this:
$$x_1(t) = a_1 \cos( \omega t + \phi_1 )$$
$$x_1 = a_1 \cos( \omega t ) \cos(\phi_1)- a_1 \sin( \omega t ) \sin(\phi_1)$$
The $$\omega$$s will be the same, so:
$$x_2 = a_2 \cos( \omega t ) \cos(\phi_2)- a_2 \sin( \omega t ) \sin(\phi_2)$$
You can now rearrange and add them together:
\begin{align} x_3 &= x_1 + x_2 \\ &= \left[a_1 \cos(\phi_1) + a_2 \cos(\phi_2) \right] \cos( \omega t ) - \left[a_1 \sin(\phi_1) + a_2 \sin(\phi_2) \right] \sin( \omega t ) \end{align}
Since you know the $$\phi$$s and $$a$$s, you can put this in the form:
$$x_3 = C \cos( \omega t ) - D \sin( \omega t )$$
Now, you just have to reverse the process to get it back into time phase form.
$$\theta = \operatorname{atan2}(D,C)$$
$$M = \sqrt{ C^2 + D^2 }$$
Therefore:
$$x_3 = M \cos( \omega t + \theta )$$
$$x_3 = M \cos\left( \omega \left( t + \frac{\theta}{\omega} \right) \right)$$
Now, just plug and chug and you should have it.
The latter part might be better understood in a different order:
\begin{align} x_3(t) &= a_3 \cos( \omega t + \phi_3 )\\ &= a_3 \cos( \omega t ) \cos(\phi_3)- a_3 \sin( \omega t ) \sin(\phi_3) \end{align}
From the equations above:
\begin{align} C &= a_3 \cos(\phi_3)\\ D &= a_3 \sin(\phi_3) \end{align}
From there it follows:
$$a_3 = \sqrt{C^2+D^2}$$
and
$$\phi_3 = \tan^{-1}\left(\frac{D}{C}\right)$$
The thing is, both of these last two equations have two possible solutions each and they need to be matched. Conveniently, the atan2 function is available on most platforms which yields the correct angle from the latter for the positive root of the former.
This solution is actually the proof of a very important principle: When two pure tones of the same frequency are added together the result is a pure tone of the same frequency although it is possible for it to have zero amplitude.
In order for the zero amplitude case to occur, the two tones have to have the same amplitude and be a half cycle ($$\pi$$ radians, 180 degrees) out of phase. This is called complete destructive interference.
If the two tones are in phase, the amplitudes are merely added and the phase remains the same. This is technically called complete constructive interference.
By repeated application of this principle, the same holds true for any linear combination of tones.
This shows up in a surprising number of places.
a. Perhaps the discussion at: https://www.dsprelated.com/showarticle/635.php would be of some value to you.
• You beat me to it, I remember reading this article a few years ago! – Ben Sep 28 '19 at 16:49
• Welcome back to this forum. Looking forward to seeing what you have to say in the next item in your answer (the one that will begin "b.") when you get around to adding it. – Dilip Sarwate Sep 28 '19 at 17:03
• Hi Dilip. (Hope you are doing well.) I don't know what you mean by "the next item in your answer (the one that will begin "b.")." – Richard Lyons Oct 2 '19 at 10:49
• @DilipSarwate, Hey Rick, you want to start your comment as I did if you want the little notification thingy to work. Your post begins: "a. Perhaps..". I, too, have been breathlessly awaiting what might be the "b." – Cedron Dawg Oct 2 '19 at 18:15
• @CedronDog, Hi Cedron. Ah, ...now I see. I don't remember typing the two 'a.' characters in my Answer. So my eyes didn't actually see them in my Answer. Thanks for helping me. In any case Mbaz can obtain his desired results by using the top row of Table 1 in my blog at my above posted web link, or he can use your equations. Happily, your and my equations produce identical results. – Richard Lyons Oct 3 '19 at 10:34 | 2020-02-18T13:01:20 | {
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https://math.stackexchange.com/questions/3476881/finding-cosine-when-negating-its-angle-given-only-the-value-of-sine | # Finding cosine when negating its angle, given only the value of sine.
If $$\sin(\theta) = \frac{1}{\sqrt{5}}$$, I must find $$\cos(-\theta)$$.
Using a right triangle (hypotenuse is $$\sqrt{5}$$ and opposite is $$1$$), it is easy to find that the adjacent side is $$2$$. Then since $$\cos(-\theta) = \cos(\theta)$$ I can simply find the answer as $$\frac{2}{\sqrt{5}}$$.
This is correct according to my homework. However, I was thinking, what if $$\theta$$ is in the second quadrant of the unit circle? Basically $$\pi > \theta > \pi/2$$.
Then, visually speaking if you look at the unit circle, there is no way that $$\cos(\theta)$$ can be positive. It has to be negative because the angle ends up in the third quadrant when you negate it as $$-\theta$$.
But my homework says that the only answer is the positive version, so I'm not sure what am I doing wrong.
• If you are only given that $\sin(\theta)=1/\sqrt 5$, there are two possible values for $\cos(\theta)$. It is possible that the answers are incomplete and/or the teacher did not think about negative values. – D.R. Dec 15 '19 at 5:10
I think you have put more good thought into this than the author of your homework answers did, and I agree with your conclusion.
Representing the angle $$\theta$$ on the unit circle, we have one end of the angle at $$(1,0)$$ and the other end at any point where $$y=\frac1{\sqrt5}$$, which is to say the other end can be at either $$\left(\frac2{\sqrt5},\frac1{\sqrt5}\right)$$ or $$\left(-\frac2{\sqrt5},\frac1{\sqrt5}\right)$$.
There are infinitely many angles with sine $$\frac1{\sqrt5}$$, since you can go around the circle as many times as you like in either direction to reach one of those two points, but you must end at one of those two points.
Then the other end of the angle $$-\theta$$ is at $$\left(\frac2{\sqrt5},-\frac1{\sqrt5}\right)$$ or $$\left(-\frac2{\sqrt5},-\frac1{\sqrt5}\right)$$.
In the end there are two possible values for $$\cos(-\theta)$$ found by taking the $$x$$ coordinates of the last two points: $$\frac2{\sqrt5}$$ and $$-\frac2{\sqrt5}.$$
(Or you could realize that $$\cos(-\theta)=\cos(\theta)$$ and therefore when you found the endpoints of the angle $$\theta$$ you already had the answer.)
But if someone did not think about it as much as you did, they might not think about the answer $$-\frac2{\sqrt5}.$$ I guess that's what happened to the author of your homework answers.
A big reason why I was so explicit about using the unit circle in this answer is that it encourages thinking about all the quadrants, whereas the right-triangle definition of sine and cosine really only makes sense in the first quadrant. | 2020-01-17T19:39:04 | {
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https://www.quizover.com/course/section/decomposing-a-composite-function-into-its-component-by-openstax | # 3.4 Composition of functions (Page 6/9)
Page 6 / 9
## Finding the domain of a composite function involving radicals
Find the domain of
Because we cannot take the square root of a negative number, the domain of $\text{\hspace{0.17em}}g\text{\hspace{0.17em}}$ is $\text{\hspace{0.17em}}\left(-\infty ,3\right].\text{\hspace{0.17em}}$ Now we check the domain of the composite function
The domain of this function is $\text{\hspace{0.17em}}\left(-\infty ,5\right].\text{\hspace{0.17em}}$ To find the domain of $\text{\hspace{0.17em}}f\circ g,\text{\hspace{0.17em}}$ we ask ourselves if there are any further restrictions offered by the domain of the composite function. The answer is no, since $\text{\hspace{0.17em}}\left(-\infty ,3\right]\text{\hspace{0.17em}}$ is a proper subset of the domain of $\text{\hspace{0.17em}}f\circ g.\text{\hspace{0.17em}}$ This means the domain of $\text{\hspace{0.17em}}f\circ g\text{\hspace{0.17em}}$ is the same as the domain of $\text{\hspace{0.17em}}g,\text{\hspace{0.17em}}$ namely, $\text{\hspace{0.17em}}\left(-\infty ,3\right].$
Find the domain of
$\left[-4,0\right)\cup \left(0,\infty \right)$
## Decomposing a composite function into its component functions
In some cases, it is necessary to decompose a complicated function. In other words, we can write it as a composition of two simpler functions. There may be more than one way to decompose a composite function , so we may choose the decomposition that appears to be most expedient.
## Decomposing a function
Write $\text{\hspace{0.17em}}f\left(x\right)=\sqrt{5-{x}^{2}}\text{\hspace{0.17em}}$ as the composition of two functions.
We are looking for two functions, $\text{\hspace{0.17em}}g\text{\hspace{0.17em}}$ and $\text{\hspace{0.17em}}h,\text{\hspace{0.17em}}$ so $\text{\hspace{0.17em}}f\left(x\right)=g\left(h\left(x\right)\right).\text{\hspace{0.17em}}$ To do this, we look for a function inside a function in the formula for $\text{\hspace{0.17em}}f\left(x\right).\text{\hspace{0.17em}}$ As one possibility, we might notice that the expression $\text{\hspace{0.17em}}5-{x}^{2}\text{\hspace{0.17em}}$ is the inside of the square root. We could then decompose the function as
We can check our answer by recomposing the functions.
$g\left(h\left(x\right)\right)=g\left(5-{x}^{2}\right)=\sqrt{5-{x}^{2}}$
Write $\text{\hspace{0.17em}}f\left(x\right)=\frac{4}{3-\sqrt{4+{x}^{2}}}\text{\hspace{0.17em}}$ as the composition of two functions.
$\begin{array}{l}g\left(x\right)=\sqrt{4+{x}^{2}}\\ h\left(x\right)=\frac{4}{3-x}\\ f=h\circ g\end{array}$
Access these online resources for additional instruction and practice with composite functions.
## Key equation
Composite function $\left(f\circ g\right)\left(x\right)=f\left(g\left(x\right)\right)$
## Key concepts
• We can perform algebraic operations on functions. See [link] .
• When functions are combined, the output of the first (inner) function becomes the input of the second (outer) function.
• The function produced by combining two functions is a composite function. See [link] and [link] .
• The order of function composition must be considered when interpreting the meaning of composite functions. See [link] .
• A composite function can be evaluated by evaluating the inner function using the given input value and then evaluating the outer function taking as its input the output of the inner function.
• A composite function can be evaluated from a table. See [link] .
• A composite function can be evaluated from a graph. See [link] .
• A composite function can be evaluated from a formula. See [link] .
• The domain of a composite function consists of those inputs in the domain of the inner function that correspond to outputs of the inner function that are in the domain of the outer function. See [link] and [link] .
• Just as functions can be combined to form a composite function, composite functions can be decomposed into simpler functions.
• Functions can often be decomposed in more than one way. See [link] .
## Verbal
How does one find the domain of the quotient of two functions, $\text{\hspace{0.17em}}\frac{f}{g}?\text{\hspace{0.17em}}$
Find the numbers that make the function in the denominator $\text{\hspace{0.17em}}g\text{\hspace{0.17em}}$ equal to zero, and check for any other domain restrictions on $\text{\hspace{0.17em}}f\text{\hspace{0.17em}}$ and $\text{\hspace{0.17em}}g,\text{\hspace{0.17em}}$ such as an even-indexed root or zeros in the denominator.
why {2kπ} union {kπ}={kπ}?
why is {2kπ} union {kπ}={kπ}? when k belong to integer
Huy
if 9 sin theta + 40 cos theta = 41,prove that:41 cos theta = 41
what is complex numbers
give me treganamentry question
Solve 2cos x + 3sin x = 0.5
madras university algebra questions papers first year B. SC. maths
Hey
Rightspect
hi
chesky
Give me algebra questions
Rightspect
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Vandna
What does this mean
cos(x+iy)=cos alpha+isinalpha prove that: sin⁴x=sin²alpha
cos(x+iy)=cos aplha+i sinalpha prove that: sinh⁴y=sin²alpha
rajan
cos(x+iy)=cos aplha+i sinalpha prove that: sinh⁴y=sin²alpha
rajan
is there any case that you can have a polynomials with a degree of four?
victor
***sscc.edu/home/jdavidso/math/catalog/polynomials/fourth/fourth.html
Oliver
can you solve it step b step
give me some important question in tregnamentry
Anshuman
what is linear equation with one unknown 2x+5=3
-4
Joel
x=-4
Joel
x=-1
Joan
I was wrong. I didn't move all constants to the right of the equation.
Joel
x=-1
Cristian
y=x+1
gary
what is the VA Ha D R X int Y int of f(x) =x²+4x+4/x+2 f(x) =x³-1/x-1
can I get help with this?
Wayne
Are they two separate problems or are the two functions a system?
Wilson
Also, is the first x squared in "x+4x+4"
Wilson
x^2+4x+4?
Wilson
thank you
Wilson
Wilson
f(x)=x square-root 2 +2x+1 how to solve this value
Wilson
what is algebra
The product of two is 32. Find a function that represents the sum of their squares.
Paul
if theta =30degree so COS2 theta = 1- 10 square theta upon 1 + tan squared theta
how to compute this 1. g(1-x) 2. f(x-2) 3. g (-x-/5) 4. f (x)- g (x)
hi
John
hi
Grace
what sup friend
John
not much For functions, there are two conditions for a function to be the inverse function: 1--- g(f(x)) = x for all x in the domain of f 2---f(g(x)) = x for all x in the domain of g Notice in both cases you will get back to the element that you started with, namely, x.
Grace | 2018-09-22T22:29:09 | {
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https://math.stackexchange.com/questions/3524292/how-is-the-inductive-hypothesis-in-strong-mathematical-induction-different-from | # How is the inductive hypothesis in strong mathematical induction different from that in ordinary induction?
I don't understand how supposing that $$P(k), k\geq 1$$ for ordinary induction is different from $$P(i), 1 \leq i \leq k, k\geq1$$ for strong induction. Example from quora:
Let’s say you wanted to prove that every positive integer has a prime factorization $$𝑝_1𝑝_2𝑝_3...𝑝_𝑚$$.
Let 𝑃(𝑛) be the statement that an integer 𝑛 has a prime factorization. We’ll proceed by strong induction. The basis is pretty clear, so I’ll leave it out.
Next we’ll assume that 𝑃(1),𝑃(2),𝑃(3),...,𝑃(𝑘) are true. 𝑘+1 can either be prime or composite, and if it’s prime we’re done, so we’ll assume it’s composite. That means 𝑘+1 can be written as a product of two positive integers, i.e. 𝑘+1=𝑝𝑞, with $$𝑝,𝑞∈ℤ^+$$. We can write 1<𝑝<𝑘+1, and 1<𝑞<𝑘+1, which implies that 2≤𝑝≤𝑘 and 2≤𝑞≤𝑘.
Here is why we need strong induction: if we had simply supposed 𝑃(𝑛) was true for arbitrary 𝑛, we would be stuck. However, we supposed that 𝑃(𝑛) was true for every positive integer up to 𝑛=𝑘, so we have much more information to work with. Because we supposed this, we know that 𝑃(𝑝) and 𝑃(𝑞) are true, i.e. that 𝑝 and 𝑞 can be represented as a product of primes. We were able to reduce the problem down to a point where 𝑝 and 𝑞 were in a range, and since our inductive hypothesis in strong induction supposes that 𝑃(𝑛) is true for a range of values (rather than just one arbitrary 𝑛), we can now use it to prove the truth of 𝑃(𝑘+1).
Using ordinary induction, I'd say that $$P(p)$$ and $$P(q)$$ are true because $$2≤𝑝≤𝑘$$ and $$2≤𝑞≤𝑘$$ and $$P(k), k\geq 1$$. Why can I not use ordinary induction here?
Another example is the proof that McCarthy 91 function equals 91 for all positive integers less than or equal to 101. The property is $$P(n)=M(101-n), n \geq 0$$ and $$M(n)$$ is the McCarthy function. The author of the proof calculates the base case for $$P(0)$$, then does a supposition that $$P(i), 0 \leq i \leq k, k \geq 0$$. The use of strong induction is justified by the fact that we need the inductive hypothesis to hold for $$k-10$$, but I don't see why $$P(n), n\geq0$$ wouldn't hold for $$n=k-10, k\geq11$$, that is $$n$$ is at least 1, if ordinary induction was used.
• For ordinary induction, you wouldn't assume $P(p)$ and $P(q)$. When proving the theorem for $k+1$, you would assume $P(k)$, and only $P(k)$ (otherwise you are doing strong induction again). And that doesn't help you at all for the prime numbers. – Dirk Jan 27 '20 at 11:40
• @Dirk Technically, you can do it with weak induction. Weak and strong are probably equivalent. However, in this case, using weak looks like a complete mess. – Arthur Jan 27 '20 at 11:45
• @Dirk right, in weak induction I make an assumption that the property holds only for the term immediately before (k+1) and prove that it holds for (k+1). p and q from the proof above do not necessary equal k, therefore I can't assume that P(q) and P(p) are true, and that's where strong induction helps with its "extended" assumption. Is it correct? – super.t Jan 27 '20 at 14:25
Ordinary or weak induction proves $$Q(n)$$ for all $$n\ge1$$ with a base step $$n=1$$ and an inductive step from $$n=k$$ to $$n=k+1$$.
Complete or strong induction considers the special case where $$Q(n)$$ denotes "$$P(k)$$ for all $$k$$ from $$1$$ to $$n-1$$ inclusive". If we try to prove $$Q(n)$$ for all $$n\ge1$$ by weak induction, the base step is vacuously true, and the inductive step is showing that, if $$P(k)$$ for all $$k$$ from $$1$$ to $$n-1$$ inclusive, then $$P(n)$$. If we can prove this statement, the weak induction on $$Q$$ succeeds, and we have also proven $$P(n)$$ for all $$n\ge1$$.
In other words, strong induction states this: if "$$P(k)$$ for all $$k$$ from $$1$$ to $$n-1$$" implies $$P(n)$$, then $$P(n)$$ for all $$n\ge1$$. Usually, the $$n-1$$ is called $$n$$ instead, so we need to prove "$$P(k)$$ for all $$k$$ from $$1$$ to $$n$$" implies $$P(n+1)$$.
Unlike weak induction, strong induction does not in general need a base step. However, in some cases the argument proving its inductive step has to consider small values of $$n$$ as special cases. | 2021-04-20T20:59:32 | {
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https://math.stackexchange.com/questions/1974188/surjections-using-stirling-numbers | # Surjections using Stirling-numbers
I have a task in which I can't really see what I'm doing wrong. The recommended answer differs from mine and it kind of makes sense to me, but I can't see what's wrong with my solution. The problem is as follows;
In a family of 4 all the 7 deadly sins are present. Every family member exercises(?)
atleast 1 of the 7 sins, but there are no sins that two distinct family members 'exercise'.
(The original problem formulation is not in english, sorry)
The task is to find out how many combinations there are if no single family member can exercise both greed and gluttony.
Basically what I did was to find out the total amount of combinations without taking the constraint into consideration (Surjection from a 7-subset to a 4-subset; $4! * S(7,4)$. After that I wanted to subtract the 'illegal' combinations. The way I figured I would find them out is as follows:
There are 4 possible family members that can exercise both greed and gluttony. After that we have 5 additional sins to 'hand out' amongst all the remaining people (including the one that we already 'handed' both gluttony and greed) so I figure'd we'd then get the total amount of 'illegal' combinations as $4 * 4! * S(5,4)$ but in the recommended solution the amount of illegal combinations was $4! * S(6,4)$ - which I understand (I guess you 'group up' both gluttony and greed as one sin and count the surjections).
But I don't understand why my solution is wrong. Any help greatly appreciated, sorry for the wall of text!
• I don't think so ("But there are no sins that two distinc family members 'exercise') - shouldn't that point to the contrary? Maybe I translated it bad but it says in the text that one specific sin can't be 'exercised' by two family members :) – Nyfiken Gul Oct 18 '16 at 14:36
• Oops; wasn’t thinking clearly. But I do now see the problem: $S(5,4)$ counts only the partitions of the other five sins into four non-empty pieces, so your calculation misses the combinations in which one person has greed and gluttony and nothing else. – Brian M. Scott Oct 18 '16 at 14:38
• Aha! That actually makes perfect sense, didn't consider that at all. Thanks again Brian! ;-) I'll gladly mark it answered if you post your comment as an answer, you're the man! – Nyfiken Gul Oct 18 '16 at 14:44
The problem here is that $S(5,4)$ counts only partitions of the other $5$ sins into $4$ non-empty parts; when you then distribute the greed/gluttony pair to one of those parts, you ensure that the family member who gets that pair also gets at least one other sin. Thus, you’re missing all of the arrangements in which one family member gets greed and gluttony and nothing else. | 2019-09-23T13:58:15 | {
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https://math.stackexchange.com/questions/4246494/finding-the-gradient-of-the-restricted-function-in-terms-of-the-gradient-of-the/4258463#4258463 | # Finding the gradient of the restricted function in terms of the gradient of the original function
The following question showed up as part of a proof that I am doing for my research thesis.
If we have a differentiable function $$f: \mathbb{R}^n \to \mathbb{R}$$ and then set $$n-d$$ coordinates to zero we get a new differentiable function $$g: \mathbb{R}^d \to \mathbb{R}$$. Now, given the gradient $$\nabla_x f(x)$$, how one can get $$\nabla_y g(y)$$?
My try
Let $$x \in \mathbb{R}^n$$ and $$S \subset \{1,\dots,n\}$$ such that $$|S|=d$$ where $$|\cdot|$$ is the cardinality of the set. Let $$U_S$$ be a restricted identity matrix such that the $$j$$-th entry of the diagonal matrix is maintained if $$j \in S$$ otherwise it is set to zero. Also, let $$I_S$$ be the restriction of $$U_S$$ where we keep nonzero columns and remove zero columns. Hence,
$$g(y)=f(U_Sx)$$ where $$y=I_S^{\top}x$$.
The above is the translation of what I stated in terms of functions $$f$$ and $$g$$.
From this point things are a little bit unclear. I think the answer should be $$\nabla_y g(y)=I_S^{\top} \nabla_x f(x)$$ but I do not know how to get it.
Also, I know using the chain rule $$J_x f(U_S x)=J_{W} f(W)J_x W= J_{W} f(W)U_S$$ where $$J$$ is the Jacobian and $$W=U_S x$$. In addition, $$\nabla^{\top}_x f(U_Sx) = J_x f(U_S x)=J_{W} f(W)U_S$$. I do not know how to put things together.
Since no one has posted an answer yet, and I get the same result as you suggest, I thought I'll post my solution for you to judge:
We have that $$U_S x = I_S y$$ so that, vieweing matrices as linear transformations $$g(y) = f(U_S x) = f(I_S y) = f\circ I_S (y)$$ And similar to what you write about $$J_{U_S}(x)$$ we have $$J_{I_S}(y) = I_S$$. Applying the chain rule: $$J_{h_1 \circ h_2}(a) = J_{h_1}(h_2(a))J_{h_2}(a)$$ then gives \begin{align} (\nabla_y g(y))^T = J_g(y) =\\ J_{f\circ I_S}(y) = \\ J_f(I_s y)J_{I_S}(y) = \\ J_f(U_s x)I_S = \\ (\nabla_x f(U_S x))^TI_S \implies \\ \nabla_y g(y) = [(\nabla_x f(U_S x))^TI_S]^T = I_S^T\nabla_x f(U_S x) \end{align}
Due to the definitions of $$U_S$$ and $$I_S$$, the zero columns in $$I_S^T$$ exactly matches the rows where $$\nabla_x f(U_S x)$$ and $$\nabla_x f(x)$$ might differ, so finally we obtain $$\nabla_y g(y) = I_S^T\nabla_x f(U_S x) = I_S^T \nabla_x f(x)$$
Edit:
As pointed out in the comments, it would be more correct to write $$\nabla_y g(y) = I_S^T\nabla_x f(I_S y)$$
• Could you make it a little more clear what's going on? Sep 17, 2021 at 10:49
• @Mathemagician314 My answer was a bit confused, I at least had some unnecessary steps. Are there any steps in particular you find strange? Sep 17, 2021 at 11:04
• @Paradox: it is not correct since $U_S\in \mathbb{R}^{n \times n}$ so $U_Sx$ is a vector in $\mathbb{R}^n$ and $I_sx \in \mathbb{R}^{d}$ Sep 18, 2021 at 16:48
• @Sepide I assume you mean $I_Sy \in \mathbb{R}^d$ since I've not written $I_S x$ anywhere. As far as I can see, $I_s$ is an $n \times d$ matrix, not $d \times n$, since it was the non-zero columns that was removed from $U_S$, not non-zero rows. So $I_s y \in \mathbb{R}^n$. Sep 18, 2021 at 16:56
Let fat matrix $${\bf S} \in \Bbb R^{d \times n}$$ be
$${\bf S} := \begin{bmatrix} {\bf I}_d & {\bf O} \end{bmatrix} {\bf P}$$
where $${\bf P}$$ is an $$n \times n$$ permutation matrix. Note that
$${\bf S} {\bf S}^\top = \begin{bmatrix} {\bf I}_d & {\bf O} \end{bmatrix} \underbrace{\,{\bf P} {\bf P}^\top}_{= {\bf I}_n} \begin{bmatrix} {\bf I}_d \\ {\bf O} \end{bmatrix} = {\bf I}_d$$
Let vector fields $$\rho : \Bbb R^n \to \Bbb R^d$$ and $$\eta : \Bbb R^d \to \Bbb R^n$$ be defined by
$$\rho := ({\bf x} \mapsto {\bf S} {\bf x}), \qquad \eta := ({\bf y} \mapsto {\bf S}^\top {\bf y})$$
and note that $$\rho \circ \eta = \mbox{id}_{\Bbb R^d}$$. Colloquially, if one "expands" and then "restricts", one ends up exactly where one started.
Given differentiable scalar field $$f : \Bbb R^n \to \Bbb R$$, let scalar field $$g : \Bbb R^d \to \Bbb R$$ be defined by
$$g := f \circ \eta$$
Hence,
\begin{aligned} g \left( {\bf y} + {\rm d} {\bf y} \right) = f \left( {\bf S}^\top {\bf y} + {\bf S}^\top {\rm d} {\bf y} \right) &= f \left( {\bf S}^\top {\bf y} \right) + \left\langle \nabla f \left( {\bf S}^\top {\bf y} \right) , {\bf S}^\top {\rm d} {\bf y} \right\rangle \\ &= f \left( {\bf S}^\top {\bf y} \right) + \left\langle {\bf S} \, \nabla f \left( {\bf S}^\top {\bf y} \right) , {\rm d} {\bf y} \right\rangle \end{aligned}
and, thus, the gradient of $$g$$ is
$$\nabla g \left( {\bf y} \right) = \color{blue}{{\bf S} \, \nabla f \left( {\bf S}^\top {\bf y} \right)}$$
or, more succinctly,
$$\boxed{ \qquad \\ \qquad \nabla g = \color{blue}{\rho \circ \nabla f \circ \eta \qquad \\ \qquad}}$$
$$\def\p{\partial} \def\L{\left}\def\R{\right} \def\LR#1{\L(#1\R)} \def\trace#1{\operatorname{Tr}\LR{#1}} \def\qiq{\quad\implies\quad} \def\grad#1#2{\frac{\p #1}{\p #2}}$$For typing convenience, name the gradient \eqalign{ p = \grad{f}{x} \\ } and rename the matrices $$I_S\to S$$ and $$U_S\to U$$.
Also note that $$\,U=SS^T\,$$ and that $$y=S^Tx \qiq dy=S^Tdx$$ Write the differential of the function and rearrange it to recover the desired gradient. \eqalign{ g(y) &= f(Ux) \\ dg &= df \\ &= p:d(Ux) \\ &= p:\LR{U\,dx} \\ &= p:\LR{SS^T\,dx} \\ &= \LR{S^Tp}:\LR{S^T\,dx} \\ &= \LR{S^Tp}:dy \\ \grad{g}{y} &= S^Tp \\ }
In the preceding, a colon is used to denote the Frobenius product, which is a concise notation for the trace \eqalign{ A:B &= \sum_{i=1}^m\sum_{j=1}^n A_{ij}B_{ij} \;=\; \trace{A^TB} \\ A:A &= \big\|A\big\|^2_F \\ } The properties of the underlying trace function allow the terms in a Frobenius product to be rearranged in many different but equivalent ways, e.g. \eqalign{ A:B &= B:A \\ A:B &= A^T:B^T \\ C:AB &= CB^T:A = A^TC:B \\ } | 2022-06-25T16:21:10 | {
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https://math.stackexchange.com/questions/3100267/stability-of-tikhonov-regularization/3100284 | # Stability of Tikhonov Regularization
$$\underset{x\in X}{\arg\inf}\left\{||Ax-b||^2+\lambda ||x||^2\right\}$$
I have read that the solution keeps the residual $$||Ax-b||^2$$ small and is stabilized through the $$\lambda ||x||^2$$ term. Can anyone help me understand why that is? I can see that the term prevents overfitting, but I can't quite see how it helps stabilizing.
Assuming $$\|\cdot\|$$ is the $$L_2$$ norm, the solution for $$x$$ is \begin{align*} x = (A^T A + \lambda I)^{-1}A^T b \end{align*} The instability in this solution lies in the inverse. If $$A$$ have columns which are nearly linearly dependent, then $$A^TA$$ is "nearly non-invertible". In other words, the condition number will be very large. The $$\lambda I$$ helps stabilize this inverse, and will always lower the condition number.
• Yes, it is the $L^{2}$ norm that is used. That makes sense! I completely forgot that the solution was relying on $(A^{T}A+\lambda I)^{-1}$. Thank you for the response! – James Feb 4 at 20:11
In order to understand Thikonov's stabilization, it helps to first look at the ordinary least square solution $$x^*$$:
\begin{align*} x^* = (A^T A)^{-1}A^T b \end{align*}
We see that it's necessary to calculate the inverse of $$A^T A$$ and this might not be possible, if $$A$$ has nearly linearly dependent columns. But let's take a closer look onto this, by factorizing just the suspicious term by a singular value decomposition.
Then $$A^T A = U \Sigma V^T$$ where $$U$$ and $$V$$ are the eigenvectors and $$\Sigma$$ is a diagonal matrix that contains the non-zero eigenvalues. It's not of special interest here that $$U=V$$, but it is very important that the pseudoinverse $$(A^T A)^{-1}$$ is found by inverting $$\Sigma$$.
More specific, the reciprocal of each eigenvalue, let's say $$\sigma_i$$ has to be found. And this might get difficult, if two columns are nearly linear dependent. In this case, $$\sigma_i$$ is very small and the result of the division gets very large and tiny pertubations of $$\sigma_i$$ lead to large fluctuations of the inverse. It is possible to monitor such cases and as this answer already mentions the condition number is one of these indicators.
The solution that Thikonov provides to overcome the problem is simple, but very effective: just take a positive variable $$\lambda$$ and add it to the denominator. This will bound the overall result and stabilize the solution:
$$$$\Sigma_{ii}^+ = 1 / (\sigma_i + \lambda)$$$$
As we now have identified the cause of instabilities and inserted a term that prevents them, we can add the same to our known equation and role it back: $$$$U (\Sigma + \lambda I) V^T = U \Sigma V^T + \lambda U V^T = A^TA + \lambda I$$$$
And finally, we arrive at the well known:
\begin{align*} x^* = (A^T A + \lambda I)^{-1}A^T b \end{align*} | 2019-08-19T17:19:38 | {
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https://math.stackexchange.com/questions/1645532/definition-of-inverse-in-linear-and-abstract-algebra | # Definition of Inverse in Linear and Abstract Algebra
In a linear algebra text, the following is the definition of the inverse of a matrix
An $n\times n$ matrix $A$ is invertible when there exists an $n \times n$ matrix $B$ such that $$AB = BA = I_n$$
And likewise in an abstract algebra textbook, the definition of the inverse of a group is
Given that $G$ is a group with operation $*$, for each $a \in G$, there exists an element $a^{-1}$ such that $$a*a^{-1} = a^{-1}*a = e,$$ where $e$ is the identity element in $G$. Such element is called the inverse of $a$ in $G$.
Unfortunately, the second semester of abstract algebra didn't quite finalize due to low enrollment, so I'm doing independent self-study of topics I missed in Linear Algebra (as mind preparation). Here's my question:
Is it sufficient to show that $AB = I_n \;\;\implies \;\;B = A^{-1}$ and $A = B^{-1}$? Or must you check both that $AB = I_n$ and $BA = I_n$ to completely conclude that $A = B^{-1}$ and $B = A^{-1}$?
I remember on an exam, I had to prove that for a group homomorphism $\phi: G\to H$, for any $a \in G$, $\phi(a^{-1}) = [\phi(a)]^{-1}$ which I proved by asking the reader to observe that $\phi(a)\phi(a^{-1}) = \phi(aa^{-1}) = \phi(e_G) = e_H$ which because $\phi(a)\phi(a^{-1}) = e_H$, this can only mean that $\phi(a)^{-1} = \phi(a^{-1})$ by definition of inverse. And I got full points for it, but it leaves me wondering: am I supposed to check both arrangements to create the strongest possible argument?
• math.stackexchange.com/questions/3852/if-ab-i-then-ba-i – mvw Feb 8 '16 at 4:52
• For the second problem concerning the homomorphism, you already know that the inverse exits because you are in a group. Proving $\phi (a^{-1})=[\phi(a)]^{-1}$ is not equivalent to proving an inverse exists. – Oliver Jones Feb 8 '16 at 4:57
• In general, if we are following the equivalent conditions for a nonsingular matrix, since $A$ is nonsingular, we know that $A^{-1}$ exists. I suppose a better way to ask my question is: Suppose we have two elements $A$,$B$ and $A$ has an inverse; does it suffice to say that $B$ is the inverse of $A$ if we just show one arrangement (i.e., $AB = e$) gives the identity? So what it seems from other answers, the left inverse and right inverse might not be the same which is why you must check both that $AB = e$ and $BA = e$, no? – Decaf-Math Feb 8 '16 at 5:02
If you already know that $G$ is a group, then to prove that $a$ and $b$ are inverses, it is enough to check $ab = e$.
If $G$ is not a group or you don't yet know that $G$ is a group (for example, if you are trying to prove that $G$ is a group), then it is not enough to show $ab =e$. You also need to show that $ba = e$.
For matrices over a field, $AB = I$ automatically implies that $BA = I$ if $A$ and $B$ are square. Otherwise, if $A$ is $m \times n$ and $B$ is $n \times m$, where $m < n$, then $BA = I$ is impossible. This can be shown by an argument using ranks.
For square matrices over many rings, it is indeed sufficient to check that $AB=I$ to conclude $A=B^{-1}$. It is certainly true for matrices over fields like the real or complex numbers. You can read about it in the link If $AB = I$ then $BA = I$
The class of rings for which this works is called the class of stably finite rings. It is a very broad class of rings, and my guess is that you are working with such a ring.
But in the broadest generality, if you are working in some wild monoid (like a maid ring over a non-stably finite ring) the. It comes necessary to check both $ab$ and $ba$.
In any group (including groups of invertible matrices) it is sufficient to check that something is either a left inverse or a right inverse. This is because in any group, the inverse necessarily exists (by definition of something being a group) and is unique (since $ag=e$ implies $g=a^{-1}e=a^{-1}$ by multiplying on the left by $a^{-1}$, and this implies $ga=e$ by multiplying on the right by $a$).
However, in an arbitrary algebraic structure, a right inverse and a left inverse may not agree.
To show an inverse for both a (square!) matrix and a group element, you must show that it is both a left and a right inverse. Matrix multiplication and general group operations are not commutative.
For homomorphisms sending inverses to inverses, you have more information, which you employed. | 2020-01-27T09:21:27 | {
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http://math.stackexchange.com/questions/167520/basic-algebra-question | # basic algebra question
Albert repays three installments, a, b and c, of a loan that he had taken for buying an electric heater. The total of the first installment and the second installment is 150. The total of the second installment and the third installment is 200. The total of the third installment and thrice the first installment is 250. What is the third installment?
it was asked in my interview for clerical work.
-
Can you translate the given information into (three linear) equations? Could you then solve the system of equations? – David Mitra Jul 6 '12 at 15:39
So you have the system
\begin{align} a + b &= 150 \\ b + c &= 200 \\ 3a + c &= 250 \end{align}
And then you must solve it. Do you know how?
One way might be to subtract the second from the first to get $a + b - (b + c) = a - c = 150 - 200 = -50$, and then to add that to the third to get $a-c + (3a + c) = 4a= -50 + 250 = 200$. This says that $4a = 200$, so that $a = 50$. Then $b = 100$ and $c = 100$.
-
how did u decide that one should substract equation 2 from 1 ... – user286035 Jul 6 '12 at 15:49
@user286035 Because doing so eliminates the variable $c$ – ItsNotObvious Jul 6 '12 at 15:59
Systematic elimination is usually the best way to go. Sometimes there are reasonable alternatives. For example, "add" all three equations. We get $4a+2b+2c=600$. But $2b+2c=400$, so $4a=200$, and therefore $a=50$.
Now easily we see that $b=100$ and therefore $c=100$.
Remark: Or perhaps the mental arithmetic is easier if from $4a+2b+2c=600$ we conclude that $2a+b+c=300$, so $2a=100$, and therefore $a=50$.
-
i liked your way. The mental arith hint is very cool. But when to use which technique ? – user286035 Jul 8 '12 at 7:44
@user286035: For "real world" equations with usually messy coefficients, use systematic (Gaussian) elimination. In interview situations, hard to know, there likely is a trick. Exploit any symmetry. – André Nicolas Jul 8 '12 at 13:02 | 2015-05-30T15:05:32 | {
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https://math.stackexchange.com/questions/9259/find-the-sum-of-all-the-multiples-of-3-or-5-below-1000 | # Find the sum of all the multiples of 3 or 5 below 1000
How to solve this problem, I can not figure it out:
If we list all the natural numbers below 10 that are multiples of 3 or 5, we get 3, 5, 6 and 9. The sum of these multiples is 23.
Find the sum of all the multiples of 3 or 5 below 1000.
• @J.M.:Much better title thank you! – AD. Nov 7 '10 at 13:42
• @AD: Sometimes straightforward is beautiful. ;) – J. M. is a poor mathematician Nov 7 '10 at 13:44
• This is actually Project Euler problem no 1 and can be solve efficiently by using mutual inclusion exclusion. – Quixotic Nov 7 '10 at 13:45
• This is the postage stamp lemma! Every number greater than 7 can be expressed as $3x+5y$ with $x,y>0$ crazyproject.wordpress.com/2010/10/22/… – N8tron Jun 28 '12 at 13:10
The previously posted answer isn't correct. The statement of the problem is to sum the multiples of 3 and 5 below 1000, not up to and equal 1000. The correct answer is \begin{eqnarray} \sum_{k_{1} = 1}^{333} 3k_{1} + \sum_{k_{2} = 1}^{199} 5 k_{2} - \sum_{k_{3} =1}^{66} 15 k_{3} = 166833 + 99500 - 33165 = 233168, \end{eqnarray} where we have the used the identity \begin{eqnarray} \sum_{k = 1}^{n} k = \tfrac{1}{2} n(n+1). \end{eqnarray}
• The one who posted the answer 233168, please explain that answer in detail . It would be really helpful for us if you will explain. – user126936 Feb 7 '14 at 9:16
• The first two sums account for multiples of $3$ and $5$, the last sum accounts for over-counting multiples of $15$ (which can appear in either of the first two sums). – user02138 Feb 7 '14 at 13:12
• For n=10; and k =3; ( 3 * (3+1) / 2) * 3 – shiva Jul 8 '14 at 10:02
• how did you know to do that? – zurbergram Jun 30 '15 at 3:19
• There is a general principal in counting things in several sets called "inclusion-exclusion". If the things can be in any of several sets, add the totals for each set, subtract the number of things that are in exactly two of the sets, add the number of things that are in exactly three of the sets, etc. – richard1941 Sep 12 '18 at 6:10
First of all, stop thinking on the number $1000$ and turn your attention to the number $990$ instead. If you solve the problem for $990$ you just have to add $993, 995, 996$ & $999$ to it for the final answer. This sum is $(a)=3983$
Count all the #s divisible by $3$: From $3$... to $990$ there are $330$ terms. The sum is $330(990+3)/2$, so $(b)=163845$
Count all the #s divisible by $5$: From $5$... to $990$ there are $198$ terms. The sum is $198(990+5)/2$, so $(c)=98505$
Now, the GCD (greatest common divisor) of $3$ & $5$ is $1$, so the LCM (least common multiple) should be $3\times 5 = 15$.
This means every number that divides by $15$ was counted twice, and it should be done only once. Because of this, you have an extra set of numbers started with $15$ all the way to $990$ that has to be removed from (b)&(c).
Then, from $15$... to $990$ there are $66$ terms and their sum is $66(990+15)/2$, so $(d)=33165$
The answer for the problem is: $(a)+(b)+(c)-(d) = 233168$
Simple but very fun problem.
• +1 but would be a bit better with an explanation as to why one should focus on 990 instead of 1000. – jmoreno Nov 13 '16 at 7:04
• Think this answer is copied from Project Eulers' answer by Rudy: projecteuler.net/thread=1#209 – Eugene Kulabuhov Nov 25 '16 at 21:59
The multiples of 3 are 3,6,9,12,15,18,21,24,27,30,....
The multiples of 5 are 5,10,15,20,25,30,35,40,45,....
The intersection of these two sequences is 15,30,45,...
The sum of the first numbers 1+2+3+4+...+n is n(n+1)/2.
The sum of the first few multiples of k, say k+2k+3k+4k+...+nk must be kn(n+1)/2.
Now you can just put these ingredients together to solve the problem.
To find n use 1000/3 = 333 + remainder, 1000/5 = 200 + remainder, 1000/15 = 66 + remainder and then sum multiples of 3: $3\cdot 333(333+1)/2 = 166833$. multiples of 5: $5\cdot 200(200+1)/2 = 100500$ and subtract multiples of 15 $15\cdot 66(66+1)/2 = 33165$ to get 234168.
• The answer that you have provided is incorrect. Take a look at the accepted answer and see how you can improve your own post. – Jeel Shah Jul 11 '13 at 23:59
• Why can't you duplicate the intersections? – JohnOsborne Jul 29 '18 at 22:32
Well the main equation is already given above. The only question which give me trouble is that why I have to subtract the sum of 15?! Well, the answer is, 15 can be evenly divide by both 3 & 5. So the products of 15 can also be divided by those number as well! So, when you adding the numbers with Sum Of Three & Sum Of Five there are some numbers(i.e. 15,30,45,60....) which are available at both SUMMATION. So, you have to subtract at least once from the total sum to get the answer!
Hope this helps someone like me:) !!
## protected by Community♦Mar 19 '14 at 16:05
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https://www.physicsforums.com/threads/global-extrema-question.66845/ | # Global extrema question
1. Mar 11, 2005
### mathrocks
I'm supposed to find the absolute max and min for f(x,y)=sin(x) +cos(y) on the rectangle defined by 0<=x<=2*pi 0<=y<=2*pi. Can someone point me in the right direction. I got the critical points but i dont know what to do after that.
2. Mar 11, 2005
### Jimmy Snyder
This can be solved with a few elementary facts concerning the sin and cos functions and no calculus at all.
First, disregarding the rectangle for a moment, it is obvious that f(x,y) must lie between -2 and 2 inclusive for any x and y, because both sin(x) and cos(y) lie between -1 and 1 inclusive, and f(x,y) is the sum of these two.
Second, it is easy to find a pair (x,y) in the rectangle where the values of -2 and 2 are actually attained. For instance if x = PI/2 and y = 0 then f(x,y) = 2. I leave -2 for you to do.
Since f(x,y) attains the values of -2 and 2 on the rectangle, and these values cannot be exceeded for any (x,y), they must be the absolute max and min.
3. Mar 11, 2005
### mathrocks
How do I do this using calculus though, so I know what to do for other problems that arise. I know I must find the partial derviative for f(x,y) and find the critical values of each of those, which I believe are (pi/2, pi), (3*pi/2, 2*pi). Then I'm suppose to do something with the rectangle but I'm not sure what?
4. Mar 11, 2005
### Jimmy Snyder
You need to check the value of f(x,y) along the perimeter of the rectangle.
Consider the fact that the maximum value of the function f(x) = x on the interval from 0 to 1 occurs at the point x = 1 where the maximum value is f(1) = 1. At that point, f'(x) = 1, not 0. So looking for points where f'(x) = 0, while necessary, is not sufficient. Loosely speaking, you need to examine the edges of the domain as well as the points where the derivative is zero.
5. Mar 11, 2005
### mathrocks
Ok, I'm starting to understand this a little better now. Now my new question is when you have the domain be something like 0<=x<=2*pi, 0<=y<=2*pi, do you have to actually go through each value and see which combination gives you a maximum? Because in the case of f(x,y)=sin(x)+cos(y), the endpoints dont yield a maximum.
6. Mar 11, 2005
### Jimmy Snyder
In general, yes. The reason is somewhat technical, and if you don't understand the rest of this paragraph, you can just accept the answer and safely skip on to the next one. Loosely speaking again, you cannot be assured that the 'edge' of the domain allows analysis. For instance, let the domain be the interior of the rectangle along with some points along the edge thrown in for good measure. Then finding the max of a function on that domain could require you to evaluate the function at those extra points on the edge one by one. Of course, if there were a lot of them, then you could be in it for the long haul. The good news is that the teacher would have the same amount of work to do in order to check your answer, so it probably won't happen to you while you are in school.
However in particular cases you may be able to do better. In fact, for the case at hand, you can look for the extrema along each of the 4 edges by looking for points where the partials are 0 and then evaluate the function at each of the 4 corners.
7. Mar 11, 2005
### mathrocks
That's the thing, when I look for the points where the partials are 0, I get
(x=pi/2), (x=3*pi/2) and (y=pi), (y=2*pi). And when I evaluate the function at the corners which are (x=0), (x=2*pi), (y=0), (y=2*pi), I get x=0, x=0, y=1, y=1. So I dont see how I can get -2,2 even though I know I should.
8. Mar 11, 2005
### Jimmy Snyder
So you need to look at the following points:
1. The points in the interior of the rectangle where the partials are 0.
2. The points along the sides where the partials are 0.
3. The corners.
The largest value of f(x,y) among the union of these three sets of points is the absolute maximum. Similarly for the minimum value.
9. Mar 11, 2005
### dextercioby
$$f(x)=:\sin x+\cos y$$
The condition for extremum-------->critical points.
$$\frac{\partial f}{\partial x} =0 ;\frac{\partial f}{\partial y}= 0$$
The type of extremum---------->
Hess(f)|_{sol.critical points} ? 0
Daniel.
P.S.If the hessian is the positive,then it's a minimum,if it's negative,it's a maximum,if it's zero,it's a saddle point... | 2018-01-22T16:40:58 | {
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https://math.stackexchange.com/questions/16366/linear-algebra-cube-dimensions-3/3360754 | # Linear Algebra, cube & dimensions > 3
I have found interesting problem in Gilbert's Strang book, ,,Introduction to Linear Algebra'' (3rd edition):
How many corners does a cube have in 4 dimensions? How many faces? How many edges? A typical corner is $(0,0,1,0)$
I have found the answer for corners:
We know, that corner is $(x_1,x_2,x_3,x_4)$. For every $x_i$ we can use either $1$ or $0$. We can do this in $2 \cdot 2 \cdot 2 \cdot 2 = 2^4 = 16$ ways.
The same method can be used for general problem of cube in $n$ dimensions (I suppose):
Let's say, we have $n$-dimensional cube (I assume, that length of edge is $1$, but it can be some $a$, where $a \in \mathbb{R}$ [1]). Here, corner of this cube looks like this: $(x_1,x_2, \ldots , x_n)$. For every $x_i$ there are $2$ possibilities: $x_i = 0$ or $x_i = 1$ ($x_i = a$ in general). So, this cube has $2^n$ corners.
It was pretty simple, I think. But now, there are also faces and edges. To be honest, I do not know, how to find the answer in this cases.
I know, that solution for this problem is:
A four-dimensional cube has $2^4 = 16$ corners and $2 \cdot 4 = 8$ three-dimensional sides and $24$ two-dimensional faces and $32$ on-dimensional edges.
Could You somehow explain me, how to figure out this solution? I have found solution for corners by myself, using Linear Algebra methods & language. Could You show me, how to find the number of edges and faces, using Linear Algebra methods?
Is there other method to find these numbers? (I suppose, that answer for this question is positive)
I am also interested in articles/textbooks/etc. about space dimensions, if You know some interesting positions about that, share with me (and community).
As I wrote: I am interested in mathematical explanations (in particular using Linear Algebra methods/language but other methods may be also interesting) and some intuitions (how to find solution using imagination etc. [2]).
Thank You for help.
[1] I am not sure of this assumption, because:
(a) I am not sure, how edges (and faces) behave in $n$ dimensions
(b) I am not sure, how should I think about the distance in $n$ dimensions. I mean, I know, that my intuition may play tricks here
[2] I am not asking, how to imagine $4$ dimensional cube, but I think, that there is a way to find the solution, using reasoning, not only Linear Algebra.
My definition of face (there was a comment about that) is the same as definition here: http://en.wikipedia.org/wiki/Face_(geometry), especially:
In geometry, a face of a polyhedron is any of the polygons that make up its boundaries.
• You are right that cubes have the same shape (number of faces) regardless of the side length. The usual distance is $\sqrt{(x_1-y_1)^2+(x_2-y_2)^2 \ldots +(x_n-y_n)^2}$ The article that Joseph Malkevitch cites is a good combinatorial approach. – Ross Millikan Jan 6 '11 at 3:10
• @Ross Millikan: I know, that using this formula I can find length in $n$ dimensions. But I can't figure out, how to imagine this, like 1,2,3 dimensions... – exTyn Jan 6 '11 at 15:30
• Gilbert Strang has his lectures on line for free at: ocw.mit.edu/courses/mathematics/… – Tpofofn Jan 7 '11 at 3:26
As you mentioned, a vertex (or a 0-face) is just a choice of string $(x_1,\ldots, x_n)$, where $x_i\in \{0,1\}$. If you think back to dimensions 2 or 3 (or even 1!), you'll see that an edge (i.e., a 1-face) is determined by two vertices which differ in a single slot. So e.g. on the unit square, the left edge goes between $(0,0)$ and $(0,1)$ -- we may name this edge more succinctly as $(0,*)$. This method generalizes, too: the bottom square (2-face) of the 3-cube is $(*,*,0)$, etc. Note that the $n$-cube will have $k$-faces for $0\leq k\leq n$. See if you can find a formula for how many there are!
This isn't exactly linear algebra, it's more like combinatorics. Also, this problem is actually a decent introduction to the idea of higher dimensions, if you try to visualize things in 4 dimensions. Search the internet for some representations of 4-cubes, and try to understand why they look the way they do.
To calculate the number of edges: as you say there are $2^n$ corners. Each one is connected to n other corners. Adding all of these up counts each edge twice, so there are $2^{n-1}n$ edges, which equals $32$ for $n=4$. Another way to count edges is to define $E(n)$ as the number of edges in $n$ dimensions. If you think of the $n+1$ dimensional cube as connecting the corresponding corners of two $n$ dimensional cubes, the recurrence is $E(n+1)=2E(n)+2^n$
Square faces are made by starting with a square in the one of the two $n$ cubes plus the translation of the edges of the $n$ cube in the new dimension. So $S(n+1)=2S(n)+E(n)$.
$3$-cube faces are made by starting with a $3$-cube in one of the two n-cubes plus translation of the squares in the new dimension. So $C(n+1)=2C(n)+S(n)$.
• An expository article about how to think about n-cubes from a combinatorial point of view can be found here: york.cuny.edu/~malk/tidbits/n-cube-tidbit.html – Joseph Malkevitch Jan 5 '11 at 1:23
• @Joseph: Thanks. It shows that the recurrence will be the same for all dimensions of boundary: points, lines, squares, etc. – Ross Millikan Jan 5 '11 at 3:36
First you need to be clear about what you mean by face, edge, vertex, etc. In my view if we have an $n$ dimensional cube, then $n-1$ dimensional cubes form its faces, and $n-2$ dimensional cubes form its edges. Vertices as you define above are 0-D points. We do not have names for $n-3$, $n-4$ dimensional boundaries so I will not attempt to name them here.
If you are looking for an intuitive way to build up the relations you can define inductive relationships using the method of extrusion. You can start with $n=1$ if you like. In this case we have one 1-D cube (i.e. a segment) with two 0-D faces (i.e. points). So
$$C_1 = 1$$ $$F_1 = 2$$ $$E_1 = 0$$ $$V_1 = 2$$
Where $C_1$ is the number of cubes in 1-D space, $F_1$ is the number of faces, $E_1$ is the number of edges (not defined here), and $V_1$ is the number of vertices.
We create a 2-D cube (i.e. a square) by extruding this 1-D cube in a direction orthogonal all current dimensions. By doing this two things happen:
1) Every object is copied maintaining its current dimensions. We had one segment, now we have two, we had two vertices, we now have four.
2) Every object extruded into the new direction creates one object of one more dimension. The 1-D segment is extruded over a second dimension to form a square (i.e. 2-D object). The vertices (0-D) are extruded to become lines.
So the number of cubes, faces, edges, vertices becomes
$$C_2 = 1$$ $$F_2 = 2C_1 + E_1 = 4$$ $$V_2 = 2V_1 = 4$$ $$E_2 = V_2$$
Repeat the process for a 3-D cube. Again we have two sources for new objects, copying and extrusion.
$$C_3 = 1$$ $$F_3 = 2C_2 + E_2 = 6$$ $$E_3 = 2F_2 + V_2 = 12$$ $$V_3 = 2V_2$$
Continue the process to get to higher dimensions.
This is the 27th problem in Exercises 1.1 of the book, Introduction to Linear Algebra, Gilbert Strang. I hope my answer will be helpful for future readers.
Question 1: How many corners does a cube have in 4 dimensions?
Solution: A corner of a cube in $$N$$ dimenison(s) can be written as $$(x_1, x_2, \cdots, x_N)$$ and for each $$x$$, the choice is either $$1$$ or $$0$$. Then the number of corners equal to $$2^N$$. Therefore, the number of corners in 4 dimesions is $$16$$.
Question 2: How many edges does a cube have in 4 dimensions?
Solution:
Notice that there are two useful invariants.
• Each edge is shared by exactly $$2$$ corners.
• Each corner is connected to $$N$$ corners.
Then the number of edges of a cube in $$N$$ dimensions is $$2^N \cdot N \cdot \frac{1}{2} = 2^{N-1}\cdot N$$. Therefore, the number of edges is $$2^{4-1} \cdot 4 = 32$$.
Question 3: How many 3D faces does a cube have in 4 dimensions?
Solution: Let's first consider the question, "How man 2D faces does a cube have in 3 dimensions?" This answer is $$3 * 2 = 6$$. $$3$$ means 3 dimensions and $$2$$ means positive and negative directions.
Now, by a similar reasoning, the number of 3D faces is $$4*2=8$$.
Question 4: How many 2D faces does a cube have in 4 dimensions?
Solutions:
Notice that there are two useful invariants.
• Starting from a corner, the number of faces between edges is $$C_{N}^{2}$$.
• Each face is shared by $$4$$ corners.
Then the numebr of 2D faces of a cube in $$N$$ dimensions is $$C_{N}^{2} \cdot 2^{N} \cdot \frac{1}{4} = C_{N}^{2} \cdot 2^{N-2}$$. Therefore, the number of 2d faces of a cube in 4 dimensions is $$6*16/4=24$$. | 2020-04-01T00:20:21 | {
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https://math.stackexchange.com/questions/1028731/if-x-subseteq-a-cup-b-then-x-subseteq-a-or-x-subseteq-b | If $X \subseteq A \cup B$, then $X \subseteq A$ or $X \subseteq B$.
If $X \subseteq A \cup B$, then $X \subseteq A$ or $X \subseteq B$.
My counterexample: Let $A = \{1\}$ and $B = \{2\}$. Then $\{1, 2\} \subseteq A\cup B$, but $\{1,2\} \not\subseteq A$ and $\{1,2\} \not\subseteq B$. How would I prove this generally? I've tried starting with the fact that if $X \subseteq A \cup B$ then $x \in X$ implies $x \in A \cup B$ and further $x \in A$ or $x \in B$, but I'm not making any decent progress after that.
How would you prove this without a counterexample?
• You already proved it: the claim is false, and it is enough to give one single counter example...and you did it. – Timbuc Nov 19 '14 at 4:42
• I know that I already proved it. I just want to do try doing it without the counterexample. – St Vincent Nov 19 '14 at 4:44
• To disprove something it suffices to find a counterexample. – Empiricist Nov 19 '14 at 4:46
• @StVincent, you can't. First, because sometimes it is true that $\;X\subset A\;\;or\;\; X\subset B\;$, and second because you must give a counter example to refute a mathematical claim. – Timbuc Nov 19 '14 at 4:47
• I know that the counterexample is sufficient, but I wanted some more practice. In other words, could it be shown false without the counterexample? – St Vincent Nov 19 '14 at 4:48
As my comment has already mentioned, counterexamples are enough to disprove a statement. If you want to characterize when the statement is wrong, an answer goes as follows:
The statement is in general wrong, given that $B \setminus A$ and $A \setminus B$ are non-empty.
Take $X = A \cup B$. Then $X \subseteq A \cup B$ trivially.
Since $B \setminus A$ is non-empty, pick $b \in B \setminus A$. Then $b \in X \setminus A$ and therefore it is false to say $X \subseteq A$.
It is symmetric to show that it is false to say $X \subseteq B$ either. This completes the proof of the statement "Given $B \setminus A$ and $A \setminus B$ are non-empty, $X \subseteq A \cup B$ does NOT imply $X \subseteq A$ or $X \subseteq B$". But please note that we also suggest a counterexample here.
Moreover, this is the largest generality you can get, i.e., the conditions "$B \setminus A$ and $A \setminus B$ are non-empty" cannot be dropped. If, say, $B \setminus A$ is empty, then $B \subseteq A$ and $A \cup B = A$. Then certainly $X \subseteq A \cup B$ implies $X \subseteq A$.
Your counterexample proves that the statement is false.
Counterexamples are perfectly valid counterproofs.
There is nothing else required--you have shown that it is not the case that the claim is true in all cases. | 2019-08-25T00:23:32 | {
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http://mathhelpforum.com/statistics/202020-sum-probabilities-combinations.html | # Math Help - Sum of Probabilities of combinations
1. ## Sum of Probabilities of combinations
Hi,
I'm thinking if you can help me think of a better way to solve this.
Here's the situation
I have probabilities of the number of bags I can sell in a day
2 - 30 %
3 - 28%
4 - 20%
5 - 16%
6 - 6%
Now, I want to know the probability of selling, for example, 8 bags after 3 days.
What I did so far is this,
I get all the permutations of 3 numbers out of {2, 3, 4, 5, 6}, where the sum is 8, so I have the following with their respective probabilities
2,2,4 = .30*.30*.20 = .018
2,3,3 = .30*.28*.28 = .02352
2,4,2 = .30*.20*.30 = .018
3,2,3 = .28*.30*.28 = .02352
3,3,2 = .28*.28*.30 = .02352
4,2,2 = .20*.30*.30 = .018
thus, the probability of selling 8 bags after 3 days is .018+.02352+.018+.02352+.02352+.018 = .12456
My problem is, I need to get all the possible number of bags sold in 3 days with their probabilities of occurring. Is there a way that I can compute for that easily? The true situation is working with more than 5 possible number of bags in a day.
2. ## Re: Sum of Probabilities of combinations
Originally Posted by rikari
Hi,
I'm thinking if you can help me think of a better way to solve this.
Here's the situation
I have probabilities of the number of bags I can sell in a day
2 - 30 %
3 - 28%
4 - 20%
5 - 16%
6 - 6%
Now, I want to know the probability of selling, for example, 8 bags after 3 days.
What I did so far is this,
I get all the permutations of 3 numbers out of {2, 3, 4, 5, 6}, where the sum is 8, so I have the following with their respective probabilities
2,2,4 = .30*.30*.20 = .018
2,3,3 = .30*.28*.28 = .02352
2,4,2 = .30*.20*.30 = .018
3,2,3 = .28*.30*.28 = .02352
3,3,2 = .28*.28*.30 = .02352
4,2,2 = .20*.30*.30 = .018
thus, the probability of selling 8 bags after 3 days is .018+.02352+.018+.02352+.02352+.018 = .12456
My problem is, I need to get all the possible number of bags sold in 3 days with their probabilities of occurring. Is there a way that I can compute for that easily? The true situation is working with more than 5 possible number of bags in a day.
It depends on your definition of "easy". The easiest way I know of to solve this problem is to use a generating function. In this case, the generating function is
$f(x) = 0.3 x^2 + 0.28 x^3 + 0.2 x^4 + 0.16 x^5 + 0.06 x^6$
I think you can see how the function is derived from your data. Then if you want to know the probability of selling exactly 8 bags in 3 days, the steps are
1. Expand $(f(x))^3$. This is easy if you have a computer algebra program, or you can just use Wolfram alpha.
2. Extract the coefficient of $x^8$ in the result. You will see in this case, that the coefficient is 0.12456, as you said.
If you would like to see the complete result of the expansion, go to Wolfram alpha
Wolfram|Alpha: Computational Knowledge Engine
and enter
expand (0.3 x^2 + 0.28 x^3 + 0.2 x^4 + 0.16 x^5 + 0.06 x^6)^3
3. ## Re: Sum of Probabilities of combinations
Thanks for the response.
I didn't think of using the generating functions, haven't thought of it yet, but surely a bit easy than what I'm currently doing.
A huge thanks | 2016-07-24T07:02:09 | {
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http://mathhelpforum.com/pre-calculus/52345-finding-slope-line.html | # Math Help - Finding the slope of a line
1. ## Finding the slope of a line
We just started doing inequalities in class and I get the basics, but this is what gets me.
M = 7, point (0,-5)
How would I write this in Y = mx + b format?
2. $y = m * x + b$
First, we know m:
$y = 7 * x + b$
Next, we know x and y:
$-5 = 7 * 0 + b$
Solve for b:
$-5 = b$
Notice, if $x = 0$, and $(x,y)$ is on the line, $b = y$.
So, now, write your equation down:
$y = 7 * x + -5$
Or:
$y = 7 * x - 5$
3. Originally Posted by largebabies
We just started doing inequalities in class and I get the basics, but this is what gets me.
M = 7, point (0,-5)
How would I write this in Y = mx + b format?
The y intercept = -5, therefore
y=7x-5
4. m=7, (0,-5)
y- (-5)
_______ = 7 : then cross multiply
x-0
y+5= 7x: move five to the other side by subtracting it
so then you get y=7x -5
5. Thank you all, it's help alot
Although, I need help on another equation, what if there is two points?
example: M = 3/2, point (5,-6)?
6. $y = m * x + b$
First, we know m:
$y = \frac 3 2 * x + b$
Next, we know x and y:
$-6 = \frac 3 2 * 5 + b$
Solve for b:
$-6 - \frac {15} 2 = \frac {-27} 2 = b$
So, now, write your equation down:
$y = \frac 3 2 * x + \frac {-27} 2$
7. Thank you once again, although it's supposed to be 27/2 not -27/2
"passes through (4,2) and (0,8)"
9. Originally Posted by largebabies
"passes through (4,2) and (0,8)"
Find the slope:
$m=\frac{y_1-y_2}{x_2-x_1}=\frac{8-2}{0-4}=\frac{6}{-4}=-\frac{3}{2}$
From (0, 8) we know the y-intersept = 8.
Use the slope-intercept form of the general equation:
$y=mx+b$, where m is the slope, and b is the y-intercept.
$y=-\frac{3}{2}x+8$
10. that helped alot!
although i'm now stuck on another equation , I did the same thing with this and got y= -10/-2 when I needed -5/4.
The equation goes like this.
“passes through (-5,10) and (3,0)”
help? | 2015-04-18T19:54:41 | {
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http://math.stackexchange.com/questions/848762/formulation-and-computation-of-the-unique-median-of-an-even-sized-list | # Formulation and computation of “the” unique median of an even-sized list
Consider an even-sized set of numbers $X = \{x_k\}$, such as $X = \{1, 2, 7, 10\}$.
The median $m$ is defined as:
$$m = \mathrm{arg \min_x} \sum_k \lvert x_k - x\rvert^1$$
Any $m \in [2, 7]$ is a minimizer of this function, and is therefore "a" median of this list.
Now, it is common practice to take the average of 2 and 7 and call it "the" median.
But that's lame, and I think I have invented (?) a more logical way to find a unique median $m^*$:
$$m^* = \lim_{\epsilon \to 0^+} \mathrm{arg \min_x} \sum_k \left\lvert x_k - x\right\rvert^{1+\epsilon}$$
Diferentiation to find the minimum only gets us so far:
$$\sum_k \mathrm{sgn}{\left(x_k - m^*\right)}\left\lvert x_k - m^*\right\rvert^{\epsilon} = 0$$ This expression can be solved numerically for smaller and smaller $\epsilon$ to give $m^* \approx 4.85$ in this example, and I suspect the "correct" median is in fact $m^* = 34/7$, but I don't know how to prove it.
I have 3 questions:
1. First of all, is this a well-known and/or useful approach? Does it have a name?
I came up with the new formulation myself, but I've never seen it used anywhere.
2. Is there some way to directly find the exact value of $m^*$, without numerical optimization?
If not, is there a better/faster approach than brute-force numerical optimization techniques?
3. Is this a (convex?) optimization problem, and if not, can it be reformulated as one?
The trouble here is that I can't find any objective function that has a unique minimum at $m^*$.
The best I can do is to find a generalized function (i.e., the limit of another function), but when I do that, I don't think the problem is a convex optimization problem anymore.
Is there another way to pose the problem that conforms better to existing optimization frameworks?
-
I am confused. If $m^*$ is the minimizer of a function over the elements $x\in X$, then how could $m^*=4.85\notin X$? – NotNotLogical Jun 26 at 20:36
@NotNotLogical: Where did you get the idea that we must have $x \in X$? – Mehrdad Jun 26 at 20:39
Ah, I was misreading. I think I now understand. – NotNotLogical Jun 26 at 20:39
What's the problem with non-uniqueness here? – Peter Sheldrick Jun 26 at 21:48
@PeterSheldrick: Huh? The problem is the very fact that it's not unique, even though there's a perfectly sensible definition that I demonstrated gives a unique answer. It doesn't make sense to say 2 is a median of that set, nor does it make sense to say 7 is, because neither is in the "center" of the data, even though they all minimize the typical objective. Hence it makes sense to look for a better number and thus a better objective. – Mehrdad Jun 26 at 21:52
If the values are sorted $x_1 \le x_2 \ldots \le x_n$, then the value $m^*$ is the unique solution on the interval $[x_{n/2},x_{n/2+1}]$ to the following equation:
$$(m^* - x_1)(m^*-x_2)\ldots(m^*-x_{n/2})=(x_{n/2+1}-m^*)(x_{n/2+2}-m^*)\ldots(x_{n}-m^*).$$
When $n=2$, it's just the mean, and when $n=4$, $m^* = (x_3x_4-x_1x_2)/(x_3+x_4-x_1-x_2)$, which in your example is $(7\cdot 10-1\cdot2)/(7+10-1-2)=34/7$. I don't see any simple way to solve the equation in closed form for higher $n$, other than standard techniques for finding roots of polynomials.
To prove that the above equation defines $m^*$, you just need to go few steps further in the manipulation of the derivative. That is $m^* = \lim_{\epsilon\rightarrow 0^+} m_\epsilon$, where $m_\epsilon$ is the solution to:
$$\sum_{k=1}^n \mathrm{sgn}{\left(x_k - m_\epsilon\right)}\left\lvert x_k - m_\epsilon\right\rvert^{\epsilon} = 0$$
For small $\epsilon$, we should have $x_{n/2}\le m_\epsilon \le x_{n/2+1}$, so this becomes:
$$-\sum_{k=1}^{n/2} (m_\epsilon-x_k)^{\epsilon} + \sum_{k=n/2+1}^{n} (x_k-m_\epsilon)^{\epsilon}=0$$
Expanding to first order in $\epsilon$ gives: $$-\sum_{k=1}^{n/2} (1+\epsilon \log(m_\epsilon-x_k)) + \sum_{k=n/2+1}^{n} (1+\epsilon\log(x_k-m_\epsilon))=O(\epsilon^2)$$
The constant terms cancel, and dividing by $\epsilon$ gives: $$-\sum_{k=1}^{n/2} \log(m_\epsilon-x_k) + \sum_{k=n/2+1}^{n} \log(x_k-m_\epsilon)=O(\epsilon)$$
Then by taking the limit as $\epsilon \rightarrow 0^+$, we get: $$-\sum_{k=1}^{n/2} \log(m^*-x_k) + \sum_{k=n/2+1}^{n} \log(x_k-m^*)=0$$ Which is equivalent to the stated condition.
There are different ways to approximate the 1-norm with a differentiable function, and each approximation will give a different unique "median". I don't know of any reason to prefer any one approximation over another other than convenience.
-
+1 Holy cow, this looks exactly like the kind of answer I was hoping for! So you used approximated $m^* - x_k$ to first-order to be equal to its linearization $1 + \epsilon \log(m^* - x_k)$, because it's equal in the infinitesimal case? It seems so obvious in hindsight but it's very clever, I wouldn't have thought of it for quite a long time! Thanks so much, I learned something new today from your answer. :) – Mehrdad Jun 28 at 0:30
Glad to help. I rewrote the argument to be more rigorous. Hopefully it's clearer now. – p.s. Jun 28 at 1:18
Indeed, very nice! Well done. – Michael Grant Jun 28 at 3:36
I'll bet that this is very amenable to a simple numerical search (Newton and/or bisection) with the midpoint as an initial condition. – Michael Grant Jun 28 at 3:39
It also makes me wonder why first-order approximations are so special. A second-order approximation would seem correct but unhelpful, whereas a zeroth-order approximation would tell us nothing. So what's so special about first-order that gives us the answer we want exactly in the limiting case? Maybe I should ask that as a question... – Mehrdad Jun 28 at 11:54
Unfortunately, this approach is not compatible with convex optimization in practice.
The reason is that in an optimization context, a convex function and its epigraph are assumed interchangeable. That is to say: consider the following two problems: $$\begin{array}{ll} \text{minimize} & f(x) \end{array}$$ $$\begin{array}{ll} \text{minimize} & y \\ \text{subject to} & f(x) \leq y \end{array}$$ These problems are equivalent if $f$ is convex: that is, given the solution to one, the solution to the other is evident, and vice-versa. Of course, the second one has an associated dual variable while the first one does not, but that doesn't change the equivalence.
Now let's consider your function for $f$, set in the second form above: $$\begin{array}{ll} \text{minimize} & y \\ \text{subject to} & \lim_{\epsilon\rightarrow 0^+} \sum_k | x_k - x |^{1+\epsilon} \leq y \end{array}$$ This is, in all practical respects, equivalent to $$\begin{array}{ll} \text{minimize} & y \\ \text{subject to} & \sum_k | x_k - x | \leq y \end{array}$$ which is of course what you'd get with the standard median function. Any practical system for optimization is really not going to be able to differentiate between the two forms. You could, of course, fix $\epsilon$ to be small and nonzero, but then you've destroyed equivalence, and of course made the leap from a linear problem to a nonlinear one.
Conceptually, what is happening here is that you are preferring a particular element of the arg min set over the rest. But establishing preferences among feasible points is precisely what the purpose of an objective function is. You need to find a way to integrate your preferences more directly into your objective or constraints. For instance, if you are determine to preserve the numerical results induced by the $|\cdot|^{1+\epsilon}$ approach---in particular, if it is important that $34/7$ be the correct answer in this example---then you will not be able to use this median function in convex optimization.
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Thanks for the response, but I feel it's begging my own question when you say "You need to find a way to integrate your preferences more directly into your objective or constraints."... the entire point of my question was to figure out how I was supposed to do that, because as I had already mentioned, I already realized this was unlikely to work with standard (convex) optimization. Your answer just basically summarized the problem I was facing but didn't help me get anywhere. – Mehrdad Jun 26 at 21:06
I am certainly not claiming to have answered every bullet point. But you did ask if this could be formulated as a (convex?) optimization problem, and I answered it. – Michael Grant Jun 27 at 4:46
Besides: your definition of median is non-standard. You well know that the standard median is unique (4.5). Even if we set aside the midpoint portion of the standard statistical definition, your optimization-based definition includes two numbers (2, 7) that do not satisfy even the fundamental criteria of a median. Your definition is already a convenience, then. Certainly you are not the only one to employ this convenience in practice, but the fact remains. – Michael Grant Jun 27 at 5:08
I did not ask if this could be formulated as a convex optimization problem, I asked if it was a (perhaps convex) optimization problem, and if not, whether it could be re-formulated as such a problem. All you did was tell me, "You need to find a way to integrate your preferences more directly into your objective or constraints". In other words, you just repeated back at me the obvious fact that I need to reformulate it as a convex optimization problem, without helping me actually get anywhere. – Mehrdad Jun 27 at 5:37
I see no reason to accept my answer. It only addresses a part of your question. Nevertheless I will incorporate your edit... – Michael Grant Jun 27 at 22:59 | 2014-11-26T01:40:07 | {
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https://math.stackexchange.com/questions/2720631/applying-method-of-characteristic-equations-to-ey-u-x-u-y-u2 | # Applying method of characteristic equations to $e^y u_x + u_y = u^2$
I have the PDE $e^y u_x + u_y = u^2$, $u(x,0) = x$ for small $|y|$. Also $u = u(x(s), y(s))$.
So $\frac{dx}{ds} = e^{y}$, $\frac{dy}{ds} = 1$, and $\frac{du}{ds} = u^{2}$.
Solving the first ode: $\frac{dx}{ds} = e^{y} \to dx = e^{y}ds$ and integrate both sides to obtain $x = e^{y}s + c_1$.
The second ODE: $\frac{dy}{ds} = 1 \to dy = 1ds$ and integrate both sides to obtain $y = s + c_2$.
The third ODE: $\frac{du}{ds} = u^{2}$, the solution to this ode is $-\frac{1}{u} = s + c_3 \implies u = -\frac{1}{s+c_3}$.
Okay now I want to eliminate the parameter $s$. I solve for $s$ from the solution of the second ode: $s = y -c_2$ then substitute into the solution of the first ode: $x = e^{y}(y -c_2) + c_1 = e^{y}y - e^{y}c_2 + c_1$. Also $u = -\frac{1}{y-c_2 + c_3}$.
This is where I am stuck and am not sure what to do. I was attempting to follow https://en.wikipedia.org/wiki/Method_of_characteristics#Example but I am not doing this correctly I think.
• With the given initial condition, you can show that $c_2 = 0$. Also, you need to substitute $y$ into the characteristic equation $\frac{dx}{ds} = e^y = e^s$ and then integrate; this will give $x = e^s + c_1$ and $c_1$ is found by solving $x(0) = x_0$, where $x_0$ is the parameterisation parameter for the "initial curve". – Chee Han Apr 3 '18 at 17:40
• Why is $e^y = e^{s}$? I know $y = y(s)$. Letting $x(0) = x_0$, from the equation $x(s) = e^s + c_1$, we get $x_0 = e^0 + c_1 = 1 + c_1$. So $c_1 = x_0 - 1$. Is this right? – Taln Apr 3 '18 at 18:04
• Because $y = s + c_2$ and $y(0) = 0$ gives $c_2 = 0$. And yes, $c_1$ is correct. Now you can do the same for $\frac{du}{ds} = u^2$, with $u(0) = x_0$. Then the final step would be to eliminate $s$ and $x_0$ if possible. Depending on the question, you might also want to determine the validity of your solution. Any standard PDE books that discuss first order PDE should have more information about this. – Chee Han Apr 3 '18 at 19:43
• Even with @CheeHan's corrections, it doesn't seem like a solution exists. Are there any typos in the original PDE? – MasterYoda Apr 3 '18 at 20:48
• @MasterYoda I've added $' for small |y| '$ to my post, but other than that I see no typos. – Taln Apr 3 '18 at 20:57
## 1 Answer
$$e^yu_x+u_y=u^2$$ I agree with your equations which can be written in a summarised form : $$\frac{dx}{e^y}=\frac{dy}{1}=\frac{du}{u^2}=ds$$ A first characteristic equation comes from $\quad\frac{dx}{e^y}=\frac{dy}{1}\quad$ leading to : $$e^y-x=c_1$$ A second characteristic equation comes from $\quad\frac{dy}{1}=\frac{du}{u^2}\quad$ leading to : $$\frac{1}{u}+y=c_2$$ The general solution of the PDE, expressed on the form of implicit equation is : $$\Phi\left((e^y-x) \:,\: (\frac{1}{u}+y) \right)=0$$ $\Phi$is an arbitrary function of two variables.
Or, equivalently the general solution of the PDE on explicit form is : $$\frac{1}{u}+y=F(e^y-x)$$ where $F$ is an arbitrary function. $$u(x,y)=\frac{1}{-y+F(e^y-x)}$$
CONDITION :
$u(x,0)=x=\frac{1}{-0+F(e^0-x)}=\frac{1}{F(1-x)}$
$F(1-x)=\frac{1}{x}$
Let $\quad X=1-x\quad;\quad x=1-X\quad;\quad F(X)=\frac{1}{1-X}$
So, the function $F$ is determined.
We put it into the above general solution where $\quad X=e^y-x\quad$ then $\quad F(e^y-x)=\frac{1}{1-(e^y-x)}$
$$u(x,y)=\frac{1}{-y+\frac{1}{1-e^y+x}}$$ This is the particular solution of the PDE which fits to the boundary condition.
Note: The above calculus is consistent with your calculus. The difference is that $ds$ is eliminated at the beginning, while in your calculus $ds$ is eliminated later. This doesn't change the final result.
ADDITION, answering to the comment of MasterYoda :
Of course, I checked the solution $u(x,y)=\frac{1}{-y+\frac{1}{1-e^y+x}}$ before publishing my answer. After the MasterYoda's comment I checked it again. Definitively the solution satisfies the PDE (copy below).
• Your solution doesn't appear to satisfy the original PDE. – MasterYoda Apr 3 '18 at 21:46
• @ MasterYoda : The solution satisfy the PDE. See the addition to my main answer. – JJacquelin Apr 4 '18 at 7:26
• Thanks, looks like I did my math wrong. – MasterYoda Apr 4 '18 at 23:59 | 2020-01-26T18:07:24 | {
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https://math.stackexchange.com/questions/606924/exponential-decay-precalculus | # Exponential Decay (Precalculus)
I'm working on this problem.
The radioactive element carbon-14 has a half-life of about 5,750 years. The percentage of carbon-14 present in the remains of animal bones can be used to determine age. How old is an animal bone that has lost 40% of its carbon-14?
I'm using the formula $y = Ce^kt$ to model the decay.
To make my life easier, I consider $5,750$ years to be 1 time unit.
So knowing this is a half life, for a single time unit I can say
$$50 = 100 e^k$$
Then
$$\frac{1}{2}=e^k$$
Then log of both sides
$$\ln\frac{1}{2}=\ln e^k$$
$$\ln\frac{1}{2}= k \cdot \ln e$$
So
$$k = \ln\frac{1}{2}$$
Good so far.
Next we need to determine the $t$ after Carbon14 has decayed to $60\%$.
So
$$60 = 100 e^kt$$
which is
$$\frac{3}{5}= e^kt$$
Then I take the log of both sides
$$\ln \frac{3}{5}= \ln e^kt$$
then
$$\ln\frac{3}{5} = kt\cdot \ln e$$
simplifies to
$$\ln\frac{3}{5} = kt$$
Plugging in $k$ we get
$$t =\frac{ \ln\frac{3}{5}}{\ln\frac{1}{2}}$$
Recall that t is a unit where each unit is $5,750$ years, so the final answer is
$$5750\cdot \left(\frac{\ln\frac{3}{5}}{\ln\frac{1}{2}}\right)$$
My answer is $\approx4,237.55$ years.
The answer key says the $4,257$ years. http://www.sosmath.com/cyberexam/precalc/EA3001/EA3001.html
Where did I go wrong?
• It seems to be that site has a mistake. What you did is correct. – DonAntonio Dec 14 '13 at 20:06
• Minor comment: It would be more conventional to use $e^{-kt}$. – André Nicolas Dec 14 '13 at 20:11
## 1 Answer
I personally think you did it right. The structure and execution is correct. I checked my answer over and over again and it is correct. I think the answer key is wrong. | 2019-11-14T03:34:44 | {
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https://www.lil-help.com/questions/777/mat-142-36-conditional-probability-and-intersections | MAT 142 3.6 Conditional Probability and Intersections
# MAT 142 3.6 Conditional Probability and Intersections
398.7k points
1. The Por Mano Tile Company imports handmade tiles from Mexico. Saltillo tiles makes up 67% of their imports and the rest are talavera tiles. Based on past experience, they expect 3.7% of the Saltillo tiles to arrive broken and 5.4% of the talavera tiles to be broken.
a. Make a complete probability tree for this situation.
b. Find the probability that a tile is a talavera tile and is broken
c. Find the probability that a tile is not broken.
d. Find the probability that a Saltillo tile is broken.
e. Find the probability that a talavera tile is not broken
2. The results of the October 3rd vote on the Emergency Stabilization Act of 2008 are shown in the table below.
a. Find the probability that a congressman voted for no
b. Find the probability that a congressman was from the west and voted yes.
c. Find the probability they voted yes given that they were from the west
d. Find the probability that they were Midwestern if they voted no.
e. Find the probability that a congressman voted no if they were from the South or West.
MAT 142
407.5k points
#### Oh Snap! This Answer is Locked
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Excerpt from file: Math Tutorial MAT 142 Problem set Unit: Probability Topic: Conditional Probability and Intersections Directions: Solve the following problems. Please show your work, use proper notation, and explain your reasoning. 1. The Por Mano Tile Company imports handmade tiles from Mexico. Saltillo tiles
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#### Oh Snap! This Answer is Locked
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Excerpt from file: Week5DQ#1RegulatoryBodies Oftheseveralregulatorybodies,whichhasthemostaffectoncompanies?Why?Doboth publicandnonpublicornotforprofitorganizationscomplywiththeregulationsofall regulatorybodies?Whyorwhynot?Aretheregrayareas?Howdocompaniesassure
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Surround your text in *italics* or **bold**, to write a math equation use, for example, $x^2+2x+1=0$ or $$\beta^2-1=0$$
Use LaTeX to type formulas and markdown to format text. See example. | 2019-02-22T21:14:07 | {
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https://stats.stackexchange.com/questions/21141/bound-on-moment-generating-function | # Bound on moment generating function
This question arises from the one asked here about a bound on moment generating functions (MGFs).
Suppose $$X$$ is a bounded zero-mean random variable taking on values in $$[-\sigma, \sigma]$$ and let $$G(t) = E[e^{tX}]$$ be its MGF. From a bound used in a proof of Hoeffding's Inequality, we have that $$G(t) = E[e^{tX}] \leq e^{\sigma^2t^2/2}$$ where the right side is recognizable as the MGF of a zero-mean normal random variable with standard deviation $$\sigma$$. Now, the standard deviation of $$X$$ can be no larger than $$\sigma$$, with the maximum value occurring when $$X$$ is a discrete random variable such that $$P\{X = \sigma\} = P\{X = -\sigma\} = \frac{1}{2}$$. So, the bound referred to can be thought of as saying that the MGF of a zero-mean bounded random variable $$X$$ is bounded above by the MGF of a zero-mean normal random variable whose standard deviation equals the maximum possible standard deviation that $$X$$ can have.
My question is: is this a well-known result of independent interest that is used in places other than in the proof of Hoeffding's Inequality, and if so, is it also known to extend to random variables with nonzero means?
The result that prompts this question allows asymmetric range $$[a,b]$$ for $$X$$ with $$a < 0 < b$$ but does insist on $$E[X] = 0$$. The bound is $$G(t) \leq e^{t^2(b-a)^2/8} = e^{t^2\sigma_{max}^2/2}$$ where $$\sigma_{\max} = (b-a)/2$$ is the maximum standard deviation possible for a random variable with values restricted to $$[a,b]$$, but this maximum is not attained by zero-mean random variables unless $$b = -a$$.
• Random variables that satisfy bounds on the mgf like the one you quote are called subgaussian random variables. They play a central role, e.g., in nonasymptotic random matrix theory and some associated results in compressed sensing. See, e.g., the link in the answer here. (This obviously doesn't speak to your particular question; but, it is of a related nature.) – cardinal Jan 16 '12 at 3:59
• @cardinal, that's interesting because I sometimes see bounds that apply only to bounded random variables, and they're useful enough that people will prove things by splitting a random variable into bounded and unbounded parts; and you're telling me that all bounded random variables are subgaussian, so I wonder how many of the inequalities I've seen for bounded random variables have generalizations to subgaussian random variables. – user54038 Feb 21 '20 at 0:37
I can't answer the first part of your question, but as for extending it to random variables with nonzero means...
First, note that any r.v. $Z$ with finite range $[a+\mu,b+\mu]$ and (necessarily finite) mean $\mu$ can be transformed into an r.v. $X = Z-\mu$ that is, of course, zero mean with range $[a,b]$ (thus satisfying the conditions in your problem statement). The transformed variate has m.g.f. $\phi_X(t) = \exp\{-\mu t\}\phi_Z(t)$ (by basic properties of the m.g.f.) Multiplying both sides by $\exp\{\mu t\}$ and applying the inequality gives:
$\phi_Z(t) = \exp\{\mu t\}\phi_X(t) \leq \exp\{\mu t\}\exp\{t^2\sigma^2_\text{max}/2\} = \exp\{\mu t + t^2\sigma^2_\text{max}/2\}$
Not surprisingly, the m.g.f. of a Normal random variable with the same mean and standard deviation equal to $\sigma_\text{max}$.
Since $$e^{ t x}$$ is a convex function, by Jensen's inequality, we have $$e^{t X} \le \frac{\sigma+X}{2\sigma}e^{-t\sigma}+ \frac{\sigma-X}{2\sigma}e^{t\sigma}$$
Taking expectation of both sides of above inequality we get $$E[ e^{t X}] \le \frac{1}{2}e^{-t\sigma} + \frac{1}{2}e^{t\sigma}$$ where we used the zero-mean assumption on $$X$$. It remains to prove that $$\frac{1}{2}(e^y+e^{-y})\le e^{y^2/2}$$. (Then, replacing $$y=t\sigma$$ we arrive at $$E[ e^{t X}] \le e^{t^2\sigma^2/2}$$.)
Here is my proof for $$\quad \quad e^y+e^{-y}\le 2 e^{y^2/2} \quad \quad \text{for }y \in \mathbb{R} \quad (1)$$.
By Taylor expansion of $$e^y$$ we have $$e^y + e^{-y} = \sum_{n=0}^{\infty} \frac{y^n}{n!} + \sum_{n=0}^{\infty} \frac{(-1)^ny^n}{n!} =2\sum_{k=0}^{\infty} \frac{y^{2k}}{(2k)!}$$ $$2e^{y^2/2} = 2\sum_{k=0}^{\infty} \frac{y^{2k}}{2^k k!}$$ from which we get (1) since $$(2k)! \ge 2^k k!$$ (this can be shown by induction or Stirling bounds on k!).
س
• Welcome to CV. When you delimit $\TeX$ markup with dollar signs \$ it will be rendered nicely. I inserted them for you. It's unclear how your first inequality is an application of Jensen's inequality, though: would you mind showing it more explicitly? – whuber Feb 20 '20 at 23:57 | 2021-01-27T20:52:25 | {
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https://math.stackexchange.com/questions/2743439/number-of-possible-sequences | # Number of possible sequences
Let $(a_1,...,a_{10})$ be a sequence with $a_i \in \{1,...,10\}$ and the following properties:
$i)$ $a_1\in\{1,...,10\}$
$ii)$ $a_i\neq a_j \ \forall i,j\in \{1,...,10\}$ with $i\neq j$
$iii)$ $a_i\in \{a_1\pm 1,...,a_{i-1}\pm1\}\cap\{1,...,10\} \quad \forall i\in \{2,...,10\}$
How many sequences with these properties exist? Is it possible to generalize this for any $n$ instead of $10$?
I tried looking at the different cases for starting values: For $a_1=\{1,10\}$ there is only one possible sequence each. For $a_1=\{2,9\}$ there are 9 possible seuqences each, because you have 9 possibilities for $1$ or $10$ and the rest of the sequences is clear. But for $a_i=\{3,4,5,6,7\}$ I don't know how to go one with that way of thinking. Can someone give me a hint?
• Is it mod arithmetic; i.e., if $a_i$ is 1 is $a_{i+1}$ allowed to be 10, and if $a_i$ is 10, can $a_{i+1}$ be 1? – Mike Apr 18 '18 at 18:09
• No, that is not the case. – Tobi92sr Apr 18 '18 at 18:16
• Isn't condition i) redundant? – user251257 Apr 18 '18 at 18:17
• I think i) is saying $a_1$ has no restrictions on choice. ii) any following a_i is restricted in that it must be distinct from any of the previous iii) when combined with ii) says that all $a_1... a_k$ will be a block of consecutive letters and $a_{k+1}$ will be either be one less than the minimum or one more than the max. – fleablood Apr 18 '18 at 21:29
Some extensive hints:
1. By induction we have $$A_i := \{ a_1, \dotsc, a_i \} = \{ \min A_i, \min A_i + 1, \dotsc, \max A_i \}.$$ Thus, $a_{i+1} = \min A_i - 1$ or $\max A_i + 1$.
2. You can and must go down exactly $k:=a_1 - 1$ times. Thus, we want to know at which $k$ indices we go down. So basically we count the subsets of $\{1, \dotsc, 9\}$ with $k$ elements (which is well-known).
For the set $\{1,2,\ldots,n\}$, let $N_k(n)$ be the number of allowable sequences $(a_1,a_2,\ldots,a_n)$ such that $a_k=n$. It is easy to see that
\begin{align} N_1(n+1)&=N_1(n)\\ N_2(n+1)&=N_1(n)\\ N_3(n+1)&=N_1(n)+N_2(n)\\ N_4(n+1)&=N_1(n)+N_2(n)+N_3(n)\\ &\,\,\vdots\\ N_{n+1}(n+1)&=N_1(n)+N_2(n)+\cdots+N_n(n) \end{align}
It follows by induction that, for $n\ge2$,
\begin{align} N_1(n)=N_2(n)&=1\\ N_3(n)=1+1&=2\\ N_4(n)=1+1+2&=4\\ &\,\,\vdots\\ N_n(n)=1+1+2+4+\cdots+2^{n-3}&=2^{n-2} \end{align}
and thus the total number of allowable sequences is
$$N_1(n)+N_2(n)+\cdots+N_n(n)=1+1+2+4+\cdots+2^{n-2}=2^{n-1}$$
Let $N(k)$ be the numbers of ways to end up with a $\{a_1,...., a_k\} = \{1,2,...k\}$.
Now if you have $\{a_1,...., a_k\} = \{1,2,...k\}$ then $a_k = 1$ or $a_k = k$ (but $a_1$ may be anything). After all, you can't have group of $a_i < m$ and a group of $a_j > m$ without some $a_l = m$.
Now the numbers of ways $\{a_1,....., a_k\} = \{1,2,...k\}$ and $a_k = k$ are the number of ways $\{a_1,..., a_{k-1}\} = \{1,....,k-1\}$ and there are $N(k-1)$ ways to do that
And the numbers of ways $\{a_1,....., a_k\} = \{1,2,...k\}$ and $a_k = 1$ are the number of ways $\{a_1,..., a_{k-1}\} = \{2,....,k\}$. And there are as many ways to do that as there are ways to do $b_i = a_i -1$ so $\{b_1,....,a_{k-1}\} = \{1,....,k-1\}$. And there are $N(k-1)$ to do that.
So $N(k) = 2N(k-1)$ and as $N(1) = 1$ (as $\{a_1\} = \{1\} \iff a_1 =1$) inductively $N(k) = 2^{k-1}$ | 2019-05-20T13:16:09 | {
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https://math.stackexchange.com/questions/2473382/show-that-the-operator-has-a-chain-of-invariant-subspaces | # Show that the operator has a chain of invariant subspaces.
Let $V$ be a $n$-dimensional vector space over $\Bbb C$ and let $T:V\rightarrow V$ be any linear operator. Show that $T$ has a chain $V_0\subseteq V_1\subseteq\ldots \subseteq V_n=V$ of invariant subspaces such that $\dim V_i=i$ for $0\le i\le n$.
Here a subspace $U$ of $V$ is an invariant subspace of $V$ if $T(U)⊆U$ and in this case $T|_U:U→U$ is a linear operator on $U$. Please help me to solve this.
• Do you know the Jordan canonical form? – amsmath Oct 15 '17 at 13:48
• Yes I know it.Jordan canonical form of a matrix. But how can I proceed from there? – abcdmath Oct 15 '17 at 14:00
• I have updated my answer and it should be alright now – Guy Fsone Oct 15 '17 at 15:35
## 3 Answers
For an elementary argument not requiring the Jordan canonical form: we argue by induction on $n$. For the base case $n=0$, the statement is trivial: the required chain is $V_0 = \{ 0 \} = V$. (If you prefer to let the base case be $n=1$, then the chain there is $\{ 0 \} \subsetneq V$.)
Now, suppose $n \ge 1$. Then since $\mathbb{C}$ is algebraically complete, $T$ has at least one eigenvector; so let $x \ne 0$ be an eigenvector. Since $T x = \lambda x$ for $\lambda$ the corresponding eigenvalue, we see that $\langle x \rangle$ is an invariant subspace of $T$. Now, consider the induced operator on the quotient space, $\bar T : V / \langle x \rangle \to V / \langle x \rangle$. This is a linear operator on an $n-1$-dimensional subspace, so by inductive hypothesis, we can find a chain $V_0 \subsetneq V_1 \subsetneq \cdots \subsetneq V_{n-1} = V / \langle x \rangle$ of $\bar T$-invariant subspaces.
Now, if $\pi : V \to V / \langle x \rangle$ is the projection operator, then we conclude that $$\{ 0 \} \subsetneq \pi^{-1}(V_0) \subsetneq \pi^{-1}(V_1) \subsetneq \cdots \subsetneq \pi^{-1}(V_{n-1}) = V$$ is a chain satisfying the requirements.
• Very nice and smart! Sometimes taking the quotient space is more helpful than just taking any complementary subspace (which here might not be $T$-invariant). – amsmath Oct 15 '17 at 16:57
I assume that you know the Jordan canonical form. So, let $T$ be any linear operator in $V$ and let $J$ be its Jordan canonical form (I use the one with ones above the diagonal). Then there exists a bijective linear map $S : \mathbb C^{n\times n}\to V$ such that $T = SJS^{-1}$. I claim that $V_i = \operatorname{span}\{Se_1,Se_2,\ldots,Se_i\}$ is a chain as desired. If you know already from your lecture that the subspaces $W_i = \operatorname{span}\{e_1,\ldots,e_i\}$ form an invariant chain for $J$, then you are already done (check it!). If not, check out the following:
First, $\{V_i\}$ is obviously nested and $\dim V_i = i$ as $S$ is invertible. Concerning the invariancy, let us start with $V_1 = \operatorname{span}\{Se_1\}$. We have $TSe_1 = SJe_1 = \lambda Se_1\in V_1$, where $\lambda$ is the first eigenvalue in the Jordan form. Ok, that's settled. Now, there is either a one right to $\lambda$ in the JCF or a zero. In the second case, you have as above $TSe_2 = \mu e_2$ with $\mu$ being the second eigenvalue in the JCF (which might be $\lambda$ or not). Let us look at the first case. Then $TSe_2 = SJe_2 = S(e_1+\lambda e_2) = Se_1 + \lambda Se_2\in V_2$. So, also $V_2$ is $T$-invariant.
I hope you get the idea...
Let do this by induction on the dimension:
1. In dimension 1 it is true take $\{0\}\subset V$ with $\dim V=1$
2. Let suppose the result is true for very subspaces $W$ of dimension $\dim U<n$ and let prove that is true for $V$ of dimension $\dim = n$ .
Since we are in complex space $T$ has eigenvalues and can decompose as direct sum of eigen-spaces that is
$$V= \bigoplus_{i =1}^{p} E_{\lambda_i}\equiv E_{\lambda_1}\oplus U$$ where we suppose that $T$ that has $p>1$ eigenvalues and let $$U = \bigoplus_{i =2}^{p} E_{\lambda_i}$$ $$T =T_1\oplus T'$$ Where $T_1 =T|_{E_{\lambda_1}}$ and $T' =T|_{U}$ .
We asume $p>1:$ We recall and it is easy to show that $T( E_{\lambda_i})\subset E_{\lambda_i}$ since
$$T(\ker (T-\lambda_iI))\subset \ker (T-\lambda_iI)$$
Hence One see that $$T( E_{\lambda_1})\subset E_{\lambda_1}$$ and $$T( U)\subset U$$ since $p>1$ we have that $$r =\dim E_{\lambda_1} <n~~~and ~~~n-r =\dim U <n$$ Whence By asumption of induction,there are two chain $$\color{green}{W_0⊆W_1⊆....⊆W_r=W}$$ such that $\dim W_i = i,~~i= 0,1,\cdots ,r$ and $T_1(W_i)\subset W_i$ and $$\color{blue}{U_0⊆U_1⊆....⊆U_{n-r}=U}$$ such that $\dim U_i = i,~~i= 0,1,\cdots ,n-r$ and $T'(U_i)\subset U_i$
Now consider the chain
$$\color{blue}{W_0\oplus U_0⊆W_0\oplus U_1⊆....⊆}\color{red}{W_0\oplus U_{n-r} =U\oplus W_0} \color{green}{\subset W_1\oplus U ⊆...⊆W_r\oplus U =W\oplus U=V}$$
That is $$\begin{cases} \color{blue}{V_i~~~~~= W_0\oplus U_i}&\text{if}~~0\le i \le n-r\\ \color{green}{V_{n-r+i} = U\oplus W_i} &\text{if}~~0\le i \le r \end{cases}$$
• For $0\le i \le n-r$ we have $\color{blue}{V_i =W_0\oplus U_i }$ then $\color{blue}{T(V_i) = T_1(W_0)\oplus T'(U_i) \subset W_0\oplus U_i = V_i}$ $$\color{blue}{\dim V_i = \dim U_i +\dim W_0 = i}$$
• For $0\le i \le r$ we have $\color{green}{V_{n-r+i} =W_i\oplus U }$ then $\color{green}{T(V_i) = T_1(W_i)\oplus T'(U) \subset W_i\oplus U = V_{n-r+i}}$
and since,$\dim U=n-r$ and $\dim W_i =i$ we have $$\color{green}{\dim V_{n-r+i} = \dim U +\dim W_{i} = n-r+i}$$ - We have the chain $$V_0⊆V_1⊆....⊆V_{n}=V$$ - $T(V_i)\subset V_i$
The cases where $p= 1$ is obvious since in that cases, it means $$V= \bigoplus_{i =1}^{p} E_{\lambda_i}=E_{\lambda_i} =\ker (T-\lambda_1 I)$$
i.e $$T= \lambda_1I$$ take $$V_i=\{v_1,v_2,\cdots,v_i\}$$
Where $\{v_1,v_2,\cdots,v_n\}$ is any basis of $V$
• Sir what is $f$ here? – abcdmath Oct 15 '17 at 13:29
• But how is $\ker T^n=V$ – Learnmore Oct 15 '17 at 13:38
• Sir, do you mean that $T(Ker T)=Ker T²$. Then by your assumption $V_2⊆V_1$ which is opposite. Please elaborate it. – abcdmath Oct 15 '17 at 13:39
• You go out from a nilpotent operator. But $T$ is any operator. Also, $\dim V_i = i$ is required. – amsmath Oct 15 '17 at 13:45
• Nope. Try to do this with$$T = \begin{pmatrix}0 & 1\\0 & 0\end{pmatrix}$$. Good luck! – amsmath Oct 15 '17 at 16:48 | 2021-01-25T17:58:42 | {
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https://math.stackexchange.com/questions/2288787/local-homeomorphism | # local homeomorphism
I'm self-studying topological manifold from the begining and I have many maybe trivial questions;
I would like to show that every homeomorphism is a local homeomorphism;
A continuous map $f:X\rightarrow Y$ is said a Local Homeomorphism if every point $x\in X$ has a neighborhood $U\subset X$ such that $f(U)$ is an open subset of $Y$ and $f|_{U}:U\rightarrow f(U)$ is an homeomorphism.
I have proceeded in this way:
If $f$ is homeomorphism,$f^{-1}$ is continuous,hence for each $x \in X$ all its open neighborhoods $U_x$ are such that $f(U_x)$ is open in $Y$.
Now I need to show that $\forall x\in X$ there is an open neighbothood $U_x$ such that: $$f|_{U_x}: U_x \rightarrow f(U_x)$$$$f^{-1}|_{U_x}: f(U_x) \rightarrow U_x$$ are both continuous (in other words $f|_{U_x}$ is an homeomorphism )
I know that if a function $g:M\rightarrow N$ is continuous then every restriction $g|_A$ with $A$ open subset of $M$ is continuous; hence $\forall x\in X$ I have that $f|_{U_x}: U_x \rightarrow f(U_x)$ is continuous in every open neighbothood $U_x$ of $x \in X$; moreover $f(U_x)$ is an open subset of $Y$ and $f^{-1}$ is continuous, then $f^{-1}|_{U_x}: f(U_x) \rightarrow U_x$ is continuous too.
Could someone check if the proof is correct?
Thank you.
• Looks OK at a glance. However, you don't need to do all this work: say $f : X\to Y$ is your homeomorphism. $X$ is a neighborhood of any point $x\in X$, $f(X) = Y$ is open in $Y$, and $\left. f\right|_X = f$ is a homeomorphism. – Stahl May 20 '17 at 8:30
• @Stahl thank you a lot. Your proof is much better than mine – Simone P May 20 '17 at 8:36
• If you're happy with this, I'll post it as an answer so the question doesn't appear unanswered. – Stahl May 20 '17 at 8:36
• ok...I'm satisfied. – Simone P May 20 '17 at 8:39
Your proof looks OK at a glance. One note: you might want to make a remark about $f$ being bijective so that $f^{-1}$ actually defines a function, and then $\left.f^{-1}\right|_U$ does as well for any $U\subseteq Y$ (as $\left.f\right|_V$ will be a bijection for any $V\subseteq X$). After all, a homeomorphism is first and foremost a bijection, and you never mention/address this aspect explicitly, although you use it implicitly in your treatment of $f^{-1}$ as a function.
However, you can give a much cleaner proof: say $f:X\to Y$ is your homeomorphism. $X$ is a neighborhood of any point $x\in X$, $f(X)=Y$ is open in $Y$, and $\left.f\right|_X=f$ is a homeomorphism. | 2019-09-19T07:26:49 | {
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http://math.stackexchange.com/questions/140231/what-is-this-question-asking-and-help-me-solve-this-equation-step-by-step | # What is this question asking and Help me solve this equation step by step
Ok So I am about to take the accuplacer college level math for a college but I do not understand this problem on the practice packet. Honestly I am blank on this problem and need help step by step on how to solve it.Please help me understand it. Thanks in advance
If a ≠ b and 1/x + 1/a= 1/b , then x =
A. 1/b – 1/a
B. b – a
C. 1/ab
D. a – b/ab
E. ab/a – b
-
It is asking you to solve for $x$, given that $$\frac{1}{x} +\frac{1}{a}=\frac{1}{b}.$$
To solve for $x$, first isolate $x$ by itself on one side; for example, move that $\frac{1}{a}$ to the right. That will give you an equation of the form $$\frac{1}{x} = \text{stuff}.$$
Do the operation on the right, and then take reciprocals (or cross-multiply) to get an expression for $x$ in terms of $a$ and $b$. Then figure out which of the five options given is that expression for $x$.
You can the work as an edit to your question and we can tell you if you are doing it right or not; that will help you learn better than me doing it for you.
(For extra points, figure out exactly on which step you need to assume $a\neq b$...)
-
Clearly explained and brief still! +1. – Joe May 3 '12 at 3:52
So the equal sign (=) with a slash (/) means if a and b are "not equal" – Backtrack May 3 '12 at 4:02
@Backtrack... Ehr... yes. If you were unsure about the symbols, you should have mentioned that explicitly. – Arturo Magidin May 3 '12 at 4:04
Yes I will try this tomorrow morning and you guys can check my work. But anyway yes, the fractions are throwing me off. To isolate x should I divide 1/a by itself and also by 1/b? And then after that what? And also should I leave the letters how they are? – Backtrack May 3 '12 at 4:25
@Backtrack: To isolate $x$, you first move the $\frac{1}{a}$ to the right; This does not involve divisions of any kind, just additions/subtractions. Then you will have an expression on the right that has no $x$'s in it: just $1$s, $a$s, and $b$s. Do the operation, which will be a sum/subtraction of fraction, and write the answer as a single fraction. You should get something along the lines of $$\frac{1}{x} = \frac{\text{expression1 involving }a\text{ and }b}{\text{expression2 involving }a\text{ and }b}.$$ Then you can cross multiply to get $x$ equal to something. – Arturo Magidin May 3 '12 at 4:28
You need to solve the equation for $x$, but I imagine the fractions are throwing you off. You can get rid of the fractions by multiplying both sides by $abx$ (why? because this is the least common multiple of $a,b$, and $x$), then it's simple algebra to rearrange and solve for $x$. Do you see why it's necessary for $a\neq b$?
- | 2014-10-23T09:45:24 | {
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https://math.stackexchange.com/questions/1963793/how-many-equivalence-classes-in-the-equivalence-relation | # How many equivalence classes in the equivalence relation
Consider the equivalence relation defined on the set A = Z \ {0}, where a~b if an only if ab > 0.
I assume this means that A is the set of all integers except 0.
How many equivalence classes are there in the above equivalence relation? Describe each of the equivalence classes.
This is my first equivalence relation assignment, and I am not sure I understand equivalence relations and classes yet. But I assume that to "form" a class, I must select all pairs (a,b) that make ab = 1 (for equivalence class 1), all those that make ab = 2 (for equivalence class 2), etc.
However, would that not mean there are infinite equivalence classes since each integer number must have a class of its own?
You are reading the relation incorrectly. Write $a \sim b$, if $ab > 0$. E.g. $1 \sim 2$ since $1\cdot 2 > 0$, in fact $1 \sim a$, with any $a > 0$, since $1a > 0$. What about $a < 0$?
• Thank you for your response. I'm a bit confused about what you mean. I am forming a class for each integer greater than zero that is formed by ab. For example, the class 1 would be {(1,1), (-1,-1)} (one pair of which accounts for a < 0). Doesn't this mean there are infinite classes that can be formed since there are infinite integers greater than 0? But the problem question suggests that there are countable classes, which has me confused. – Lorraine Oct 11 '16 at 13:43
• The definition of an equivalence class is as follows: Let $X$ be a set and let $\sim$ be an equivalence relation on $X$. Then the equivalence class of $x \in X$, is $[x] = \left\{ y \in X \: \middle| \: x \sim y \right\}$. Thus, in this exercise, the equivalence class of $1$ is the set of all integers $a$ such that $1a > 0$, since $1 \sim a$ if $1a > 0$. – Matias Heikkilä Oct 11 '16 at 13:47
An equivalence relation on a set $S$ "divides" that set into disjoint subsets, called equivalence classes. Now, suppose $\sim$ is an equivalence relation and $a \in S$. The equivalence class of $a$ under $\sim$, denoted $[a]$ is defined by
$$[a] := \{ b \in S \,|\, a \sim b\}$$
So, in your example, given a non-zero integer $n$, its equivalence class is the set
$$[n]:= \{ m \in \mathbb{Z}\setminus\{0\} \, \mid \, nm > 0\}$$
By definition of an equivalence relation, $\sim$ is reflexive, which means that $a \sim a$ for every $a \in S$, so $a \in [a].$
Using this property for $2 \in \mathbb{Z}\setminus \{0\},$ we see that $2 \in [2]$, which is true since $2 \cdot 2 >0$.
$3$ is also a member of $[2]$, since $2 \cdot 3 > 0$, but $-2 \not\in [2]$, since $2 \cdot (-2) < 0,$ so we have at least two distinct equivalence classes.
Hope this helps. | 2019-09-19T09:01:40 | {
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http://mathhelpforum.com/calculus/24673-two-answers-correct.html | 1. ## Which two answers are correct?
Here's the basic question:
We have an integral, that of which we do not know, where two out of three students got the correct answer. Which two students are right?
a. $sin^2x+C$
b. $-cos^2x+C$
c. $-sin^2x+C$
My logic is that one of the $(-)sin^2x+C$ is right, meaning the integral might be $\int2sin(x)cos(x)$, but I'm not sure how one could get a $-cos^2x+C$ out of that.
Any help is appreciated.
2. An integral may have many primitives.
Why don't you differentiate the options to check if they yield the integrand.
3. Originally Posted by Krizalid
An integral may have many primitives.
Why don't you differentiate the options to check if they yield the integrand.
When you talk about primitives, do you mean a given integral may have two answers which don't actually equal each other but are nonetheless valid evaluations of the integral?
4. Hello, ebonyscythe!
This is a classic problem . . .
Integrate: . $2\int\sin x\cos x\,dx$
There are (at least) three solutions . . .
(1) Let $u = \sin x\quad\Rightarrow\quad du = \cos x\,dx$
Substitute: . $2\int u\,du \;=\;u^2 + C\;=\;\boxed{\sin^2\!x + C}$
(2) Let $u = \cos x\quad\Rightarrow\quad du = -\sin x\,dx$
Substitute: . $2\int u(-du) \:=\:-u^2 + C\:=\:\boxed{-\cos^2\!x + C}$
(3) . $\sin x\cos x \:=\:\frac{1}{2}(2\sin x\cos x) \;=\;\frac{1}{2}\sin2x$
So we have: . $2\int\frac{1}{2}\sin2x\,dx \;=\;\int\sin2x\,dx \;=\;\boxed{-\frac{1}{2}\cos2x + C}$
. . . . . and these three answers are equivalent.
5. if the choices were the answers of the 3 students and only two got the correct answer, then student who didn't get it right might had written $-\sin^2 x + C$, like what kriz wrote, differentiate the three, since we don't really know what the integral was, at least, when the three were differentiated, 2 of them would yield the same derivative..
6. Originally Posted by DivideBy0
When you talk about primitives, do you mean a given integral may have two answers which don't actually equal each other but are nonetheless valid evaluations of the integral?
yes, they only differ by constants..
7. Originally Posted by ebonyscythe
Here's the basic question:
We have an integral, that of which we do not know, where two out of three students got the correct answer. Which two students are right?
a. $sin^2x+C$
b. $-cos^2x+C$
c. $-sin^2x+C$
My logic is that one of the $(-)sin^2x+C$ is right, meaning the integral might be $\int2sin(x)cos(x)$, but I'm not sure how one could get a $-cos^2x+C$ out of that.
Any help is appreciated.
Is it not the case that:
$\sin^2(x)=1-\cos^2(x)$?
ZB
8. Originally Posted by Constatine11
Is it not the case that:
$\sin^2(x)=1-\cos^2(x)$?
ZB
it can be the case.. take note of choice B.. C takes charge the "1" in your right hand side.. | 2013-05-26T01:54:18 | {
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https://www.klarapys.com/square-mile-fmsebyy/viewtopic.php?c0f73b=how-to-find-the-degree-of-a-term | What is the degree of the following polynomial? You don't have to use Standard Form, but it helps. We can find the degree of a polynomial by identifying the highest power of the variable that occurs in the polynomial. With the help of the community we can continue to Learn how to determine the end behavior of the graph of a polynomial function. The first one is 4x 2, the second is 6x, and the third is 5. St. Louis, MO 63105. Learn how to find the degree and the leading coefficient of a polynomial expression. on or linked-to by the Website infringes your copyright, you should consider first contacting an attorney. Classification of polynomials vocabulary defined. Plus examples of polynomials. If you've found an issue with this question, please let us know. ChillingEffects.org. Quiz on Degree of Polynomial What is the degree of following polynomial? The degree of the polynomial is found by looking at the term with the highest exponent on its variable (s). If a and b are the exponents of the multiple variables in a term, then the degree of a term in the polynomial expression is given as a+b. In this non-linear system, users are free to take whatever path through the material best serves their needs. Notice our 3-term polynomial has degree 2, and the number of factors is also 2. Track your scores, create tests, and take your learning to the next level! Varsity Tutors LLC Varsity Tutors. Definition: The degree is the term with the greatest exponent. The largest degree of those is 3 (in fact two terms have a degree of 3), so the polynomial has a degree of 3, The largest degree of those is 4, so the polynomial has a degree of 4. The sum of the exponents in each term of the expansion are 3. So far what I … The term shows being raised to the seventh power, and no other in this expression is raised to anything larger than seven. When a polynomial has more than one variable, we need to find the degree by adding the exponents of each variable in each term. Here, the highest exponent is x5, so the degree is 5. an Find the degree. Example #1: 4x 2 + 6x + 5 This polynomial has three terms. With more than one variable, the degree is the sum of the exponents of the variables. (More correctly we should work out the Limit to Infinity of ln(f(x))/ln(x), but I just want to keep this simple here). The degree of a polynomial is the highest degree of its terms The leading coefficient of a polynomial is the coefficient of the leading term Any term that doesn't have a variable in it is called a "constant" term types of polynomials depends on … link to the specific question (not just the name of the question) that contains the content and a description of The degree of a polynomial is the highest degree of its terms. Step 4:The largest power of the variable is the degree of the polynomial deg(x5+x3+x2+x+x0) = 5 The given expression can be re-written as: The polynomial terms may only have variables raised to positive integer exponents. I have a question where I am asked to find the amount of terms required in a Maclaurin polynomial to estimate $\cos(1)$ to be correct to two decimal places. Baylor University, Bachelor in Arts, Philosophy and Religious Studies, General. Identify the term containing the highest power of x to find the leading term. More examples showing how to find the degree of a polynomial. There are 4 simple steps are present to find the degree of a polynomial:- Example: 6x5+8x3+3x5+3x2+4+2x+4 1. The Standard Form for writing a polynomial is to put the terms with the highest degree first. 3 x 2 + x + 33. - example: 6x5+8x3+3x5+3x2+4+2x+4 1 all Rights Reserved, SAT Courses & in! Simple steps are present to find the 7th Taylor polynomial centered at x = 0 the. 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https://www.hpmuseum.org/forum/showthread.php?mode=linear&tid=6939&pid=61973 | Simpson revival
09-28-2016, 08:04 AM
Post: #1
Pekis Member Posts: 120 Joined: Aug 2014
Simpson revival
Hello,
The problem with Simpson's rule for calculating a definite integral is that you have to specify the number of intervals. I surely reinvented the wheel but I found a way to deal with a given tolerance (e.g. 1E-5), WITHOUT wasting any calculation in the process. It's very useful when you have to call integration in the middle of a root solving problem.
Beginning with 2 intervals, it doubles the number of intervals at each step and recycles previous results.
That lead to this compact program (here in BASIC (for Namir ):
What do you think of it ?
Variables
A,B: End points
f(X): Function to integrate
P: Expected Tolerance
Z: Integral value
W: Z from previous step
N: Number of intervals
L: N from previous step
H: Size of one interval ( (B-A)/N )
S: Sum of 4*f(x) for new (odd) points for this step
R: S from previous step
U: Sum of 2*f(x) for even points, known from previous steps
T: U from previous step, initialized with f(A)+f(B)
Program
10 INPUT A
20 INPUT B
30 P=1E-5 'Expected tolerance
40 R=0: L=1: N=1: W=1E99 'Initializing
50 X=A: GOSUB 1000: T=Y 'Call Y=f(X) subroutine with X=A
60 X=B: GOSUB 1000: T=T+Y 'T=f(A)+f(B) now
70 N=N*2: H=(B-A)/N 'New step: Double the number of intervals
90 FOR I=1 TO L 'Calculate Sum of f(x) for new (odd) points
100 GOSUB 1000: S=S+Y
110 X=X+2*H 'Next new (odd) point
120 NEXT I
130 S=S*4: U=T+R/2 'Odd points have a coefficient of 4; Even points come from old odd points but with a coefficient of 2 instead of 4
140 Z=H*(S+U)/3 'New Integral value at this step
150 IF ABS(Z-W)<=P THEN PRINT "Integral value: ";Z: STOP 'OK with expected tolerance
160 W=Z: T=U: R=S: L=N 'Prepare new step
170 GOTO 70 'New step, please
1000 Y=f(X) 'Calculate f(X) subroutine
1010 RETURN
09-28-2016, 10:40 AM (This post was last modified: 09-28-2016 10:48 AM by Namir.)
Post: #2
Namir Senior Member Posts: 813 Joined: Dec 2013
RE: Simpson revival
Very clever algorithm! Hats off! The reusing of summations from previous calculations is very clever.
Namir
PS: And thank you for the BASIC code.
09-28-2016, 12:25 PM (This post was last modified: 09-28-2016 12:43 PM by Dieter.)
Post: #3
Dieter Senior Member Posts: 2,397 Joined: Dec 2013
RE: Simpson revival
(09-28-2016 08:04 AM)Pekis Wrote: The problem with Simpson's rule for calculating a definite integral is that you have to specify the number of intervals. I surely reinvented the wheel but I found a way to deal with a given tolerance (e.g. 1E-5), WITHOUT wasting any calculation in the process.
As it is the case with most clever ideas, you indeed reinvented the wheel. ;-)
Which of course doesn't make the method less clever.
(09-28-2016 08:04 AM)Pekis Wrote: What do you think of it ?
It's such an obvious improvement that some time ago I posted a HP41 version of this idea. This program even calculates an improved approximation from the two last Simpson results, using the error estimate associated with this method. For instance, if you have the results for 4 and 8 intervals, you can get a new result from these two, having an accuracy comparable to the result with 16 intervals. What about implemeting this idea in your BASIC program?
A long time ago a much better mathematician had a similar idea based on the trapezoid method combined with an extraordinarily clever extrapolation scheme, leading to what is now known as the Romberg method.
Dieter
09-28-2016, 02:00 PM
Post: #4
Namir Senior Member Posts: 813 Joined: Dec 2013
RE: Simpson revival
(09-28-2016 12:25 PM)Dieter Wrote:
(09-28-2016 08:04 AM)Pekis Wrote: The problem with Simpson's rule for calculating a definite integral is that you have to specify the number of intervals. I surely reinvented the wheel but I found a way to deal with a given tolerance (e.g. 1E-5), WITHOUT wasting any calculation in the process.
As it is the case with most clever ideas, you indeed reinvented the wheel. ;-)
Which of course doesn't make the method less clever.
(09-28-2016 08:04 AM)Pekis Wrote: What do you think of it ?
It's such an obvious improvement that some time ago I posted a HP41 version of this idea. This program even calculates an improved approximation from the two last Simpson results, using the error estimate associated with this method. For instance, if you have the results for 4 and 8 intervals, you can get a new result from these two, having an accuracy comparable to the result with 16 intervals. What about implemeting this idea in your BASIC program?
A long time ago a much better mathematician had a similar idea based on the trapezoid method combined with an extraordinarily clever extrapolation scheme, leading to what is now known as the Romberg method.
Dieter
A few years ago I replaced the Trapezoidal rule use in Romberg with Simpson's method (and also tried using other quadrature algorithms). You can check it here and click on the link A New Face of Romberg Integration . There is a link next to it that allows you to download a related Excel demo file by Graeme Dennes.
Namir
09-28-2016, 03:50 PM (This post was last modified: 09-28-2016 05:17 PM by Namir.)
Post: #5
Namir Senior Member Posts: 813 Joined: Dec 2013
RE: Simpson revival
Here is the Basic code from Pekis with the improved calculations using Diter's approach
Here is version 1 with the improvement applied only to the final results:
Code:
10 INPUT A 20 INPUT B 30 P=1E-5 'Expected tolerance 40 R=0: L=1: N=1: W=1E99 'Initializing 50 X=A: GOSUB 1000: T=Y 'Call Y=f(X) subroutine with X=A 60 X=B: GOSUB 1000: T=T+Y 'T=f(A)+f(B) now 70 N=N*2: H=(B-A)/N 'New step: Double the number of intervals 80 X=A+H: S=0 'Start with 1st new (odd) point 90 FOR I=1 TO L 'Calculate Sum of f(x) for new (odd) points 100 GOSUB 1000: S=S+Y 110 X=X+2*H 'Next new (odd) point 120 NEXT I 130 S=S*4: U=T+R/2 'Odd points have a coefficient of 4; Even points come from old odd points but with a coefficient of 2 instead of 4 140 Z=H*(S+U)/3 'New Integral value at this step 150 IF ABS(Z-W)>P THEN GOTO 160 153 PRINT "Integral value: ";Z 155 PRINT "Improved Integral: "; (16*Z-W)/15 157 STOP 'OK with expected tolerance 160 W=Z: T=U: R=S: L=N 'Prepare new step 170 GOTO 70 'New step, please 1000 Y=f(X) 'Calculate f(X) subroutine 1010 RETURN
Here is version 2 where the improvment is applied to each iteration:
Code:
10 INPUT A 20 INPUT B 30 P=1E-5 'Expected tolerance 40 R=0: L=1: N=1: W=1E99 'Initializing 50 X=A: GOSUB 1000: T=Y 'Call Y=f(X) subroutine with X=A 60 X=B: GOSUB 1000: T=T+Y 'T=f(A)+f(B) now 70 N=N*2: H=(B-A)/N 'New step: Double the number of intervals 80 X=A+H: S=0 'Start with 1st new (odd) point 90 FOR I=1 TO L 'Calculate Sum of f(x) for new (odd) points 100 GOSUB 1000: S=S+Y 110 X=X+2*H 'Next new (odd) point 120 NEXT I 130 S=S*4: U=T+R/2 'Odd points have a coefficient of 4; Even points come from old odd points but with a coefficient of 2 instead of 4 140 Z=H*(S+U)/3 'New Integral value at this step 150 IF ABS(Z-W)>P THEN GOTO 160 153 PRINT "Integral value: ";Z 155 PRINT "Improved Integral: "; (16*Z-W)/15 157 STOP 'OK with expected tolerance 160 IF W=1E99 THEN W=Z ELSE W=(16*Z-W)/15 165 T=U: R=S: L=N 'Prepare new step 170 GOTO 70 'New step, please 1000 Y=f(X) 'Calculate f(X) subroutine 1010 RETURN
Version 1 gave better results. These results seem counter intuitive, since I expected version 2 to give a better final result.
Any or all improvements and corrections are welcome.
Namir
09-28-2016, 10:00 PM (This post was last modified: 09-28-2016 10:46 PM by Dieter.)
Post: #6
Dieter Senior Member Posts: 2,397 Joined: Dec 2013
RE: Simpson revival
(09-28-2016 03:50 PM)Namir Wrote: Version 1 gave better results. These results seem counter intuitive, since I expected version 2 to give a better final result.
Any or all improvements and corrections are welcome.
I am not sure how your code works in detail, but here is an example to compare with.
Edit: I think I now know where your error is. The improved result is calculated from the previous (W) and current (Z) Simpson approximation. But in your second program W is replaced with the last improved approximation, cf. line 160. So the next improved value is calculated from the the previous improved value and the current Simpson approximation. Here seems to be the problem. You must not overwrite W. Leave it as it is and only print/display the improved value.
Or, even better, calculate the previous and current improved approximations WI and ZI and stop the iteration as soon as these agree within the specified tolerance.
Finally, here is my example:
Let f(x) = 1/x and integrate it from 1 to 2.
The exact result, rounded to 12 digits, is ln(2)=0,693147180560.
Code:
n Standard Simpson Error Improved Error ---------------------------------------------------------------------- 2 0,694444444444 1,3 E-03 4 0,693253968254 1,1 E-04 0,693174603175 2,7 E-05 8 0,693154530655 7,4 E-06 0,693147901481 7,2 E-07 16 0,693147652819 4,7 E-07 0,693147194297 1,4 E-08 32 0,693147210290 3,0 E-08 0,693147180788 2,3 E-10 64 0,693147182421 1,9 E-09 0,693147180564 3,6 E-12 128 0,693147180676 1,2 E-10 0,693147180560 5,6 E-14
The "improved" column holds the extrapolated results, i.e. [16*Simpson(a, b, n) – Simpson(a, b, n/2)] / 15.
The "Error" columns show the absolute error of the calculated integrals.
You can see that the "improved" values are roughly as accurate as the standard Simpson results with twice the number of intervals, and sometimes even better.
From step to step the error of the standard method improves by a factor approaching 16, while it's finally factor 64 for the improved method.
Here is the VBA code I used to calculate Simpson(1, 2, n):
Code:
Function Simpson(a, b, n) h = (b - a) / n k = 4 s = 0 For i = 1 To n - 1 s = s + k * f(a + i * h) k = 6 - k Next Simpson = (f(a) + s + f(b)) * h / 3 End Function Function f(x) f = 1 / x End Function
What results do you get?
Dieter
09-29-2016, 04:44 PM
Post: #7
Namir Senior Member Posts: 813 Joined: Dec 2013
RE: Simpson revival
So my version 1 of my Basic code is the way to go since it does not update the older area value in W.
To be honest, I am a bit disappointed by your algorithm, since the general expectation for improving a calculated value in an iteration should be used in the next iteration and cause enhancing the result and/or reduce the number of iteration. An example is Ostrowski's root solving algorithm where he first uses Newton's method to refine the guess for the root and then (in the same iteration) refines that guess again. The result is that Ostrokswki's method matches the efficiency of Halley's root finding method.
In version 1, the (error of the refined area/area of calculated result) is about 100! It seems that applying the refinement on two good estimates for the integral works well to enhance the final result only.
Namir
09-29-2016, 06:23 PM (This post was last modified: 09-29-2016 09:42 PM by Dieter.)
Post: #8
Dieter Senior Member Posts: 2,397 Joined: Dec 2013
RE: Simpson revival
(09-29-2016 04:44 PM)Namir Wrote: So my version 1 of my Basic code is the way to go since it does not update the older area value in W.
IMHO the way to go is the approach in my HP41 program:
You calculate the Simpson approximation for n=2, 4, 8, ... intervals. After each step the improved value is calculated from the last two approximations. This value is returned to the user. If last two improved (!) approximations agree within the stated tolerance, the iteration exits.
In the given example two "improved" steps yield a similar accuracy as three "standard" Simpson steps, so the final result is obtained faster. If only one single iteration is saved, the total execution time is reduced by 50%.
(09-29-2016 04:44 PM)Namir Wrote: To be honest, I am a bit disappointed by your algorithm, since the general expectation for improving a calculated value in an iteration should be used in the next iteration and cause enhancing the result and/or reduce the number of iteration.
Just try it the way described above. ;-)
EDIT: Here is a short VBA program that implements the suggested method.
Code:
Function Simpson(a, b, Optional errlimit = 0.000001, Optional nmax = 2000) fab = f(a) + f(b) s2 = 0 s4 = f((a + b) / 2) simp_new = (fab + 4 * s4) * (b - a) / 6 improved_new = simp_new n = 4 Do h = (b - a) / n s2 = s2 + s4 s4 = 0 For i = 1 To n - 1 Step 2 s4 = s4 + f(a + i * h) Next simp_old = simp_new simp_new = (fab + 2 * s2 + 4 * s4) * h / 3 improved_old = improved_new improved_new = (16 * simp_new - simp_old) / 15 n = 2 * n Loop Until (Abs(improved_new - improved_old) <= errlimit) Or (n > nmax) Simpson = improved_new End Function Function f(x) f = 1 / x End Function
Dieter
09-29-2016, 10:27 PM (This post was last modified: 09-29-2016 10:30 PM by Namir.)
Post: #9
Namir Senior Member Posts: 813 Joined: Dec 2013
RE: Simpson revival
A modified version of the BASIC program works if we have two tests:
If Abs(Z - W) <= P Then
....
And
If Abs(Z - (16 * Z - W) / 15) <= P Then
...
I tested f(x)=1/X for integrals in [1,2], [1,10], and [1,100]. In each case the absolute value of the ratio of the error of the calculated integral divided by the error of the improved integral is in the order of 1000!
I suggest omitting the first test altogether!
Here is an updated Basic version that reflects Dieter's suggestion and my above comments:
Code:
10 INPUT A 20 INPUT B 30 P=1E-5 'Expected tolerance 40 R=0: L=1: N=1: W=1E99 'Initializing 50 X=A: GOSUB 1000: T=Y 'Call Y=f(X) subroutine with X=A 60 X=B: GOSUB 1000: T=T+Y 'T=f(A)+f(B) now 70 N=N*2: H=(B-A)/N 'New step: Double the number of intervals 80 X=A+H: S=0 'Start with 1st new (odd) point 90 FOR I=1 TO L 'Calculate Sum of f(x) for new (odd) points 100 GOSUB 1000: S=S+Y 110 X=X+2*H 'Next new (odd) point 120 NEXT I 130 S=S*4: U=T+R/2 'Odd points have a coefficient of 4; Even points come from old odd points but with a coefficient of 2 instead of 4 140 Z=H*(S+U)/3 'New Integral value at this step 145 Y=(16*Z-W)/15 ' Calculate improved integral 150 IF ABS(Z-Y)<=P THEN PRINT "Integral value: ";Y: STOP 'OK with expected tolerance 160 W=Z: T=U: R=S: L=N 'Prepare new step 170 GOTO 70 'New step, please 1000 Y=f(X) 'Calculate f(X) subroutine 1010 RETURN
Namir
09-30-2016, 06:00 AM (This post was last modified: 09-30-2016 06:07 AM by Dieter.)
Post: #10
Dieter Senior Member Posts: 2,397 Joined: Dec 2013
RE: Simpson revival
(09-29-2016 10:27 PM)Namir Wrote: A modified version of the BASIC program works if we have two tests:
If Abs(Z - W) <= P Then
This compares the old and new Simpson result. As you already said, this test is obsolete.
(09-29-2016 10:27 PM)Namir Wrote: If Abs(Z - (16 * Z - W) / 15) <= P Then
...
This compares the last Simpson result and the last improved approximation. Does this make sense?
(09-29-2016 10:27 PM)Namir Wrote: I suggest omitting the first test altogether!
The test should compare the last and previous improved approximations, cf. the posted VBA code. In your program that's the current and previous value of Y.
Dieter
10-02-2016, 03:29 PM (This post was last modified: 10-02-2016 03:37 PM by Dieter.)
Post: #11
Dieter Senior Member Posts: 2,397 Joined: Dec 2013
RE: Simpson revival
(09-30-2016 06:00 AM)I Wrote: The test should compare the last and previous improved approximations, cf. the posted VBA code. In your program that's the current and previous value of Y.
It looks like even this remaining error can be estimated and an even more accurate result can be calculated from the last two improved values. Here is a more structured experimental version.
Code:
Function Simpson(a, b, Optional errlimit = 0.000001, Optional nmax = 3000) fab = f(a) + f(b) s2 = 0 s4 = 0 simp_new = 0 improved_new = 0 n = 2 Do h = (b - a) / n s2 = s2 + s4 s4 = 0 For i = 1 To n - 1 Step 2 s4 = s4 + f(a + i * h) Next simp_old = simp_new improved_old = improved_new simp_new = (fab + 2 * s2 + 4 * s4) * h / 3 If n = 2 Then improved_new = simp_new converged = False Else improved_new = (16 * simp_new - simp_old) / 15 converged = Abs(improved_new - improved_old) < 63 * errlimit End If n = 2 * n Loop Until converged Or n > nmax Simpson = (64 * improved_new - improved_old) / 63 End Function Function f(x) f = 1 / x End Function
What do you think?
The following diagram may give an impression for the integral of f(x) = 1/x from a=1 to b=2. Here X is the number of iterations (i.e. n=2x intervals) and Y is the accuracy, i.e. the number of valid digits. The red line is the classic Simpson method, the blue one represents the "improved" method, and finally the green graph shows the further extrapolated result as shown in the last program.
Dieter
10-02-2016, 04:48 PM
Post: #12
Namir Senior Member Posts: 813 Joined: Dec 2013
RE: Simpson revival
(10-02-2016 03:29 PM)Dieter Wrote:
(09-30-2016 06:00 AM)I Wrote: The test should compare the last and previous improved approximations, cf. the posted VBA code. In your program that's the current and previous value of Y.
It looks like even this remaining error can be estimated and an even more accurate result can be calculated from the last two improved values. Here is a more structured experimental version.
Code:
Function Simpson(a, b, Optional errlimit = 0.000001, Optional nmax = 3000) fab = f(a) + f(b) s2 = 0 s4 = 0 simp_new = 0 improved_new = 0 n = 2 Do h = (b - a) / n s2 = s2 + s4 s4 = 0 For i = 1 To n - 1 Step 2 s4 = s4 + f(a + i * h) Next simp_old = simp_new improved_old = improved_new simp_new = (fab + 2 * s2 + 4 * s4) * h / 3 If n = 2 Then improved_new = simp_new converged = False Else improved_new = (16 * simp_new - simp_old) / 15 converged = Abs(improved_new - improved_old) < 63 * errlimit End If n = 2 * n Loop Until converged Or n > nmax Simpson = (64 * improved_new - improved_old) / 63 End Function Function f(x) f = 1 / x End Function
What do you think?
The following diagram may give an impression for the integral of f(x) = 1/x from a=1 to b=2. Here X is the number of iterations (i.e. n=2x intervals) and Y is the accuracy, i.e. the number of valid digits. The red line is the classic Simpson method, the blue one represents the "improved" method, and finally the green graph shows the further extrapolated result as shown in the last program.
Dieter
You latest code works in general better than Pekis' version. The ratio of errors generated by your code to that of Pekis range in the orders [1, 10]. Good work!
10-09-2016, 03:14 AM (This post was last modified: 10-09-2016 03:15 AM by Namir.)
Post: #13
Namir Senior Member Posts: 813 Joined: Dec 2013
RE: Simpson revival
Thank you Pekis for your code. I had fun translating it into programs for the HP-71B, HP-85B, HP-41C and HP-67. I like your approach because it gives accurate results. Of course, the calculations take more time since the program repeats the integral calculations at ever decreasing increment values.
Namir
07-31-2018, 02:57 PM (This post was last modified: 07-31-2018 06:42 PM by Albert Chan.)
Post: #14
Albert Chan Senior Member Posts: 1,676 Joined: Jul 2018
RE: Simpson revival
(09-28-2016 03:50 PM)Namir Wrote: Version 1 gave better results. These results seem counter intuitive, since I expected version 2 to give a better final result.
Any or all improvements and corrections are welcome.
Hi, Namir
From the book SICP, I think I know why your version 2 does not give expected better estimate.
The trick is not just apply correction based on 2 iterations, but from all previous iterations.
(09-28-2016 10:00 PM)Dieter Wrote: Let f(x) = 1/x and integrate it from 1 to 2.
The exact result, rounded to 12 digits, is ln(2)=0,693147180560.
Code:
n Standard Simpson Error Improved Error ---------------------------------------------------------------------- 2 0,694444444444 1,3 E-03 4 0,693253968254 1,1 E-04 0,693174603175 2,7 E-05 8 0,693154530655 7,4 E-06 0,693147901481 7,2 E-07 16 0,693147652819 4,7 E-07 0,693147194297 1,4 E-08 32 0,693147210290 3,0 E-08 0,693147180788 2,3 E-10 64 0,693147182421 1,9 E-09 0,693147180564 3,6 E-12 128 0,693147180676 1,2 E-10 0,693147180560 5,6 E-14
The "improved" column holds the extrapolated results, i.e. [16*Simpson(a, b, n) – Simpson(a, b, n/2)] / 15.
The first estimate is correct: T1 + (T1-T0)/(16-1) = 0.693174603175
However, second estimate should apply correction again, from previous estimate.
T1 + (T1-T0)/(64-1) = 0.693147901481 + (-2.67017e-05)/63 = 0.693147477645
Third estimate need 3 corrections (/15, /63, /255), ... This is the revised table.
Code:
n Standard Simpson Error Improved Error ---------------------------------------------------------------------- 2 0,694444444444 1,3 E-03 4 0,693253968254 1,1 E-04 0,693174603175 2,7 E-05 8 0,693154530655 7,4 E-06 0,693147477645 3,0 E-07 16 0,693147652819 4,7 E-07 0,693147181917 1,4 E-09 32 0,693147210290 3,0 E-08 0,693147180562 2,4 E-12 64 0,693147182421 1,9 E-09 0,693147180560 1,6 E-15
Edit:
After reading the whole thread, it seems the code is not trying to use Romberg's method.
All it wanted was the first approximation of Romberg's extrapolation.
Sorry for the noise.
Anyway, this is what Romberg's method would look like ...
08-04-2018, 05:29 PM
Post: #15
Albert Chan Senior Member Posts: 1,676 Joined: Jul 2018
RE: Simpson revival
In Python, Romberg's method corrections can be simplified with generator.
Code:
from __future__ import division def simpson(f, a, b, n=2): # Simpson generator h = (b-a) / 2 t0 = (f(a)+f(b)) * h while 1: s = sum(f(a + i*h) for i in xrange(1, n, 2)) t1 = t0/2 + s*h # simpson's partial sum yield t1 + (t1-t0)/3 # simpson's formula t0, n, h = t1, 2*n, h/2
Code:
def _romberg(gen, y, k): # Romberg generator for z in gen: yield z + (z-y)/k y = z
Code:
def romberg(f, a, b, verbal=False, k=15): "Romberg's Method integration" gen = simpson(f, a, b, n=2) x = gen.next() while 1: y = gen.next() x = y + (y-x) / k if verbal: print x if x == y: return x # stop if converged gen = _romberg(gen, y, k) # Romberg corrections k = 4 * k + 3
Redo previous post example:
>>> romberg(lambda x: 1/x, 1, 2, verbal=True)
0.693174603175
0.693147477645
0.693147181917
0.693147180562
0.69314718056
0.69314718056
0.69314718055994506
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https://math.stackexchange.com/questions/2119309/fast-way-of-finding-the-remainder | # Fast Way of Finding The Remainder
I have the following question:
Find the remainder of $29\times 2901\times 2017$ divided by $17$
I already have the answer (7) for this problem. I solved it using the long way by multiplying all of the numbers then divvy them with 17. I am just thinking if there are any fast way of solving this. My solution takes a long time.
Hint:
If $$a\equiv b\pmod c$$ And $$d\equiv e\pmod c$$ You can multiply the two to get $$a\times d\equiv b\times e\pmod c$$
Doing this for the numbers separately, then multiplying the expressions will get you the solution very quickly.
• See here for this Congruence Product Rule. But here we can also exploit radix rep to mod out chunks of digits, which greatly simplifies the computation - see my answer. – Bill Dubuque Jan 29 '17 at 19:52
You can write each number as a multiple of $17$ plus a small remainder:
$$(17+12)(170\cdot 17+11)(118\cdot 17+11).$$
When you multiply this out (which you don't have to do) every term is a multiple of $17$ except the last one $12\cdot 11\cdot 11$. So you need consider only this last product. We can repeat what we just did with a little factoring. The above $= 3\cdot 11 \cdot 4 \cdot 11 = 33\cdot 44 = (17+16)(2\cdot 17 + 10).$ So you need consider only $16\cdot 10$.
• $(-5)\cdot(-6)\cdot(-6)$ is easier to compute. – tomasz Jan 29 '17 at 18:09
• @tomasz Yes, but from the question, I don't think the OP knows about modular arithmetic yet, so dealing with a negative remainder is an extra complication. Although I see a "mod" answer was accepted as best, so maybe I'm wrong. – B. Goddard Jan 29 '17 at 18:23
Find remainder after dividing $29 \times 2901 \times 2017$ by $17$.
Work $\mod 17$ !
\begin{align} & 29 \times 2901 \times 2017 \\ & 12 \times 1201 \times 317 \\ & 12 \times 1201 \times 147 \\ & 12 \times 351 \times 62 \\ & 12 \times 11 \times 11 \\ & 1452 \\ & 602 \\ & 92 \\ & 7. \end{align}
• This is much easier if one works with least magnitude reps (i.e. allow negative remainders), e.g. see my answer. – Bill Dubuque Jan 29 '17 at 19:54
This can be done in $$15$$ seconds of purely mental modular arithmetic of small numbers using the universal divisibility test, which is essentially the division algorithm ignoring quotients. To obtain optimal speedup we use least magnitude remainders, e.g. $$-1$$ vs. $$16\pmod{\!17}$$ since doing so simplifies subsequent arithmetic. To reduce a decimal number mod $$n$$ we continually mod out the leading chunks of its digits. Since we allow negative remainders, we will encounter negative digits, delimited by a comma, e.g. $$\,a,b := a(10)\!+\!b.\,$$ We prove $$\ 3247\equiv 0\pmod{\!17}\,$$ for practice.
\begin{align}{\rm mod}\ 17\!:\qquad &\,\ \color{#90f}{32}\ 47\\ \equiv\ &{\color{#90f}{-2}},\color{#0a0}47 \ \ \ \text{by }\ \ \ \ \ \,\color{#90f}{32}\,\equiv\,\color{#90f}{-2} \\ \equiv\ &\quad\ \ \, \color{#f84}{\bf 1}7\ \ \ \text{by }\ {\color{#90f}{-2}},\color{#0a0}4 \equiv\, \color{#90f}{{-}2}(10)\!+\!\color{#0a0}4\equiv -16\equiv \color{#f84}{\bf 1} \\[-.3em] \text{Let's do the number in the OP}\qquad\ \ \ \ \\[-.3em] &\,\ \color{#90f}{29}\ 01\\ \equiv\ & {\color{#90f}{-5}},\color{#0a0}01\ \ \text{ by }\quad \color{#90f}{29\,\equiv\, -5} \\ \equiv\ &\quad\ \ \color{#f84}{\bf 1}1\ \ \, \text{ by }\ \color{#90f}{{-}5},\color{#0a0}0\equiv {\color{#90f}{-5}(10)\!+\!\color{#0a0}0}\equiv -50\equiv\color{#f84}{\bf 1}\\ \equiv\ &\quad \,\color{#08f}{-6}\\[-.2em] \text{Similarly \,2017\equiv\color{#c00}{-6}\ so we have}\phantom{MM}\\[-.2em] &\ 29\cdot 2901\cdot 2017\\ \equiv\ &(-5)(\color{#08f}{-6})(\color{#c00}{-6})\\ \equiv\ &(-5)\ \color{#08f} 2\\ \equiv\ &\ 7 \end{align}\qquad\qquad
where we have applied the Congruence Product and Sum Rules many times above.
Remark I wrote every little detail above to help avoid confusion with negative digits. Once one gains proficiency there is no need to be so extremely verbose. See here for a larger example which also employs negative digits for optimization.
• (+1) The OP did say "fast way" of finding the remainder, this is it! – Silverfish Jan 29 '17 at 23:42
29 × 2901 × 2017 $\equiv$ 12 × 11 × 11 (mod 17)
12 × 121 $\equiv$ (mod 17)
12 × 2 $\equiv$ (mod 17)
24 $\equiv$ 7 (mod 17)
• Even faster if you use negative modulos to keep the numbers small, so in this case, the first line would be $(-5)*(-6)*(-6)$ – WorldSEnder Jan 29 '17 at 17:49
• Ok I keep in mind that. – Kanwaljit Singh Jan 29 '17 at 17:53
• @KanwaljitSingh: $(-a) \equiv (b - a) (\mathrm{mod}\, b)$ – Kevin Jan 29 '17 at 19:31
The numbers are given in decimal representation, therefore I would start by simplifying $100 \bmod 17$. We have that $20\equiv 3$, therefore $100=20\times 5\equiv 3\times 5 = 15\equiv -2$. This allows us to write:
$$2900\equiv -2\times 29\equiv -2\times (-5) = 10$$
And
$$2000\equiv -2\times 20\equiv-2\times 3 = -6$$
We thus have:
$$29\times 2901\times 2017\equiv -5\times 11\times (-6) = 30\times 11\equiv -4\times 11\equiv -4\times (-6) = 24\equiv 7$$ | 2021-01-25T02:27:36 | {
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https://mathhelpboards.com/threads/when-bubble-sort-to-sort-a-list-of-numbers-7-12-5-22-13-32.5715/ | # When Bubble Sort to sort a list of numbers 7, 12, 5, 22, 13, 32
#### Joystar1977
##### Active member
Can somebody tell me which example is right when a question that is given to me says to bubble sort a list of numbers 7, 12, 5, 22, 13, 32? I found two examples and one was with a graph that included Original List, Pass 1, Pass 2, Pass 3, Pass 4, Pass, 5, and Pass 6, the numbers with 7 on one line, 12 on another line, 5 on another line, 22 on another line, 13 on another line, and 32 on another line which includes Comparisons, and Swaps. The second example was to do the following:
First Pass:
(7, 12, 5, 22, 13, 32)- Here, algorithm compares the first two elements, and swaps since 7 <12.
(7, 12, 5, 22, 13, 32)- Swap since 12 > 5.
(7, 12, 5, 22, 13, 32)- Swap since 5 < 22.
(7, 12, 5, 22, 13, 32)- Swap since 22 > 13.
(7, 12, 5, 22, 13, 32)- Now, since these elements are already in order (32 > 13), algorithm does not swap them.
Second Pass and Third Pass is done the same way.
Can someone please tell me which example is correct when bubble sorting the list of numbers?
#### Joystar1977
##### Active member
Discrete Mathematics (Bubble Sort to sort the list)
Is it true that the number 32 is definitely in its correct position at the end of the first pass when it comes to the following numbers: 7, 12, 5, 22, 13, 32?
#### Joystar1977
##### Active member
Discrete Mathematics Numbers 7, 12, 5, 22, 13, 32
How does the number of comparisons required change as the pass number increases?
#### Evgeny.Makarov
##### Well-known member
MHB Math Scholar
Re: Discrete Mathematics Numbers 7, 12, 5, 22, 13, 32
Can somebody tell me which example is right when a question that is given to me says to bubble sort a list of numbers 7, 12, 5, 22, 13, 32? I found two examples and one was with a graph that included Original List, Pass 1, Pass 2, Pass 3, Pass 4, Pass, 5, and Pass 6, the numbers with 7 on one line, 12 on another line, 5 on another line, 22 on another line, 13 on another line, and 32 on another line which includes Comparisons, and Swaps. The second example was to do the following:
First Pass:
(7, 12, 5, 22, 13, 32)- Here, algorithm compares the first two elements, and swaps since 7 <12.
(7, 12, 5, 22, 13, 32)- Swap since 12 > 5.
(7, 12, 5, 22, 13, 32)- Swap since 5 < 22.
(7, 12, 5, 22, 13, 32)- Swap since 22 > 13.
(7, 12, 5, 22, 13, 32)- Now, since these elements are already in order (32 > 13), algorithm does not swap them.
Second Pass and Third Pass is done the same way.
Can someone please tell me which example is correct when bubble sorting the list of numbers?
I did not understand the description of the graph very well. The second example is incorrect. First, the swaps are not shown: all lines are the same. Second, if the sorting is done in increasing order, then there is no reason to swap 7 and 12: they are already ordered correctly.
The list during the first pass should change as follows. (Compared elements are underlined.)
\begin{array}{rrrrrr} \underline{7} &\underline{12}& 5& 22& 13& 32\\ 7& \underline{12}& \underline{5}& 22& 13& 32\\ 7& 5& \underline{12}& \underline{22}& 13& 32\\ 7& 5& 12& \underline{22}& \underline{13}& 32\\ 7& 5& 12& 13& \underline{22}& \underline{32}\\ 7& 5& 12& 13& 22& 32\end{array}
Is it true that the number 32 is definitely in its correct position at the end of the first pass when it comes to the following numbers: 7, 12, 5, 22, 13, 32?
Yes, if the sorting is done in increasing order.
How does the number of comparisons required change as the pass number increases?
In the unoptimized bubble sort, the number of comparisons is the same for all passes. In the optimized sort, the number decreases by 1 from one pass to the next because numbers at the end of the array are in their correct positions. Thus, for an array of $n$ numbers the first pass requires $n-1$ comparisons, the second $n-2$ comparisons, and so on.
#### Joystar1977
##### Active member
Thank you Evgeny.Makarov! I seen two different examples of Bubble Sorting and was trying to see which is done correctly. I wanted to just clarify which one I am supposed to use when Bubble Sorting. Am I correct then that its easier to see things on a graph for Bubble Sorting such as having the Original List of numbers with 7 on one line, 12 on another line, 5 on another line, 22 on another line, 13 on another line, and 32on another line, Pass 1, Pass, 2, Pass 3, Pass 4, Pass 5, Pass 6, Comparisons, and Swaps? Just trying to get a thorough understanding of this.
#### Evgeny.Makarov
##### Well-known member
MHB Math Scholar
Am I correct then that its easier to see things on a graph for Bubble Sorting such as having the Original List of numbers with 7 on one line, 12 on another line, 5 on another line, 22 on another line, 13 on another line, and 32on another line, Pass 1, Pass, 2, Pass 3, Pass 4, Pass 5, Pass 6, Comparisons, and Swaps? Just trying to get a thorough understanding of this.
I am having trouble understanding this description of a graph, so I can't say whether it makes sense. Execution of Bubble sort can be illustrated in many ways. The most natural, I think, is to show the array after each comparison and potential swap. I showed the first pass in this way in my previous post. For another illustration, see Wikipedia. Note that while there are many ways to illustrate the execution of the Bubble sort, the algorithm itself works in a unique, fixed way. It is this way of working that is important, not whether it is illustrated as a graph or as a sequence of lines.
#### Joystar1977
##### Active member
Thanks for the explanation Evgeny.Makarov! | 2021-06-20T19:51:31 | {
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http://math.stackexchange.com/questions/113970/evaluating-sum-n-1-infty-frac22n-15n1 | # Evaluating $\sum_{n=1}^\infty \frac{2^{2n-1}}{5^{n+1}}$
Evaluating
$$\sum_{n=1}^\infty \frac{2^{2n-1}}{5^{n+1}}$$
Its $\frac{\infty}{\infty}$ so I should use L Hospital rule, but the terms are exponential and differentiation won't do much good? I am thinking maybe I somehow use $\ln$ both sides? But how? Or perhaps I should do something else?
UPDATE
Following @Paul's answer: Since $|r| < 1$ so sequence is convergent. So I use the formula
$$\sum_{n=1}^\infty ar^{n-1} = \frac{a}{1-r}$$
$$\frac{1}{10} \cdot \frac{1}{1-\frac{4}{5}} = \frac{1}{10}\cdot 5 = \frac{1}{2}$$
But answer is $\frac{2}{5}$
-
L'Hopital's rule is for evaluating limits, not evaluating series. – JavaMan Feb 27 '12 at 13:00
What's the point of your update? If you follow Paul's answer, you'll have it. The first term is "$\frac{4}{5}$", not $1$. – Gigili Feb 27 '12 at 13:11
@Jiew Meng: Compare the series $\sum_{n=1}^\infty \left(\frac{4}{5}\right)^n=\sum_{n=1}^\infty \frac{4}{5}\left(\frac{4}{5}\right)^{n-1}$ with $\sum_{n=1}^\infty ar^{n-1}$, we have $a=\frac{4}{5}$ and $r=\frac{4}{5}$. Therefore, using the formula you have, we have $\sum_{n=1}^\infty \left(\frac{4}{5}\right)^n=\frac{\frac{4}{5}}{1-\frac{4}{5}}=4$. – Paul Feb 27 '12 at 13:16
@CliveNewstead: That's what I said! – Gigili Feb 27 '12 at 13:18
@Gigili: I know, but Jiew Meng either hadn't read or hadn't understood it, so I was reiterating. – Clive Newstead Feb 27 '12 at 14:23
You can rewrite it as $$\sum_{n=1}^\infty \frac{2^{2n-1}}{5^{n+1}}=\sum_{n=1}^\infty \frac{\frac{1}{2}\cdot2^{2n}}{5\cdot 5^n}=\frac{1}{10}\sum_{n=1}^\infty \frac{2^{2n}}{5^n}=\frac{1}{10}\sum_{n=1}^\infty \frac{4^{n}}{5^n}=\frac{1}{10}\sum_{n=1}^\infty \left(\frac{4}{5}\right)^n$$ which is a geometric series with ratio $r=\displaystyle\frac{4}{5}$. I think you can finish it from here.
-
I think there is a $\frac{1}{10}$ too in the first step. – Riccardo.Alestra Feb 27 '12 at 12:36
Actually shouldn't the $\frac{1}{10}$ appear in step 3? – Jiew Meng Feb 27 '12 at 12:44
Btw, I also updated my question – Jiew Meng Feb 27 '12 at 12:58
Sorry, that's a typo. $\frac{1}{10}$ shouldn't appear until the second equality. – Paul Feb 27 '12 at 13:13
This is a series, not a sequence. L'Hopital's rule is used for the latter, not the former.
When working with series, it is a good idea to first determine if the the given series is of a certain known type. One type is a series of the form $$\sum\limits_{n=m}^\infty a r^{n+k} =ar^{m+k} + ar^{m+k+1}++ ar^{m+k+2}+\cdots.$$
These are called geometric series and the quantity $r$ is called the ratio of the series.
A geometric series converges if and only if $|r|<1$. When the series converges, it converges to to the first term of the series divided by the quantity $(1-r)$. The first term of the series above is obtained when you take $n=m$: $ar^{m+k}$.
I would not suggest you use a formula to find the sum of the series, but rather do the following:
1. Identify the ratio $r$.
2. If $|r|<1$, the series converges to $\text{the first term of the series}\over 1-r$. If $|r|\ge1$, the series diverges.
In your case $$\sum_{n=1}^\infty { {2^{2n-1}\over 5^{n+1}} }=\sum_{n=1}^\infty\textstyle {1\over 10}({4\over5})^n$$ This series is geometric with $r=4/5$. The first term of the series is ${1\over10}\cdot{4\over5}$ (obtained by setting $n=1$). So, the series converges and $$\sum_{n=1}^\infty{\textstyle {1\over 10}({4\over5})^n} = {{1\over10}\cdot{4\over5}\over 1-{4\over5}} ={{1\over10}\cdot{4\over5}\over {1\over5}}={2\over5}.$$
Another example: $$\sum_{n=2}^\infty 4\cdot (-1/3)^{n+5}\quad \buildrel {r=-1/3}\over{ = }\quad {4\cdot(-1/3)^7\over 1-(-1/3)} ={-4/3^7\over 5/3}.$$ Here, the first term of the series is obtained when $n=2$: $4\cdot(-1/3)^{2+5}=4\cdot(-1/3)^{7}$.
And one more: $$\sum_{n=0}^\infty ( 1/3)^{n }\quad \buildrel {r= 1/3}\over{ = }\quad { ( 1/3)^0\over 1-(1/3)} ={1\over 2/3}={3\over2}.$$
-
I don't see what special points you added that Paul didn't mention in his answer or it's not mentioned in comments. – Gigili Feb 27 '12 at 13:56
@Gigili Mostly how to actually find the sum of a geometric series. It seemed to OP had problems using the formula. – David Mitra Feb 27 '12 at 13:59
@Gigli Oh, I didn't see the hidden comments... – David Mitra Feb 27 '12 at 14:00
I think the OP just made a mistake while calculating the sum, and Paul linked to the WP article about Geometric series. I once posted an answer completely different from the other answer and I recall you complained why I posted a new answer. – Gigili Feb 27 '12 at 14:04 | 2016-05-03T20:47:37 | {
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https://math.stackexchange.com/questions/25528/cubic-root-of-negative-numbers | # cubic root of negative numbers
Excuse my lack of knowledge and expertise in math,but to me it would came naturally that the cubic root of $-8$ would be $-2$ since $(-2)^3 = -8$.
But when I checked Wolfram Alpha for $\sqrt[3]{-8}$, real it tells me it doesn't exist.
I came to trust Wolfram Alpha so I thought I'd ask you guys, to explain the sense of that to me.
-8 has three cube roots: $-2$, $1 + i \sqrt{ 3 }$ and $1 - i \sqrt{ 3 }$. So you can't answer the question "Is $\sqrt[3]{-8}$ real" without specifying which of them you're talking about.
For some reason, WolframAlpha is only giving $1 + i \sqrt{ 3 }$ as an answer -- that looks like a bug in WolframAlpha to me.
• You can get all three roots from WolframAlpha with x^3=-8 but for cbrt and ^(1/3) (like sqrt and ^(1/2)) it gives a single answer which for continuity reasons is not on the negative real line. – Henry Mar 7 '11 at 15:34
• Actually, Wolfram is giving the correct principal cube root. Mathematica is much more oriented towards continuous mathematics than discrete mathematics, which makes the extension of exponentiation to odd roots of negative numbers very out of place. – Hurkyl Jan 9 '15 at 16:01
• The last bit is the only blemish in an otherwise excellent answer. – J. M. is a poor mathematician Mar 11 '16 at 13:08
Although it's been two years since this question was asked, some folks might be interested to know that this behavior has been modified in WolframAlpha. If you ask for the cube root of a negative number, it returns the real valued, negative cube root. Here, I just asked for "cbrt -8", for example:
Note "the principal root" button. That allows you to toggle back to the original behavior. Near the bottom, we still see information on all the complex roots.
We can plot functions involving the cube root and solve equations involving the cube root and it consistently acts real valued. If you just type in an equation, it will solve it, plot both sides and highlight the intersections. Here's "cbrt(x)=sin(2x)"
See this. In particular, the prinicipal cube root has nonzero imaginary part.
• I think it says the prinicipal cube root has positive imaginary part, but in practice it takes gives a non-negative real cube root of a non-negative real. In fact, looking for example at (-32)^(1/5), it takes the root with the smallest non-negative anti-clockwise rotation from the positive real axis. Looking at (-32)^(3/5) it seems to take the fifth root before cubing. – Henry Mar 7 '11 at 15:40
Of course, you're absolutely right about $-2$ being a cubic root of $-8$.
The point might be that there are actually, three different cubic roots of $-8$, namely the roots of the polynomial $x^3+8$. One of this roots is real ($-2$), the other two are complex and conjugate of each other.
If you ask Wolfram for the cubic root of $-8$ you get one of these two non real roots, namely $1+(1.732050807568877293527446341505872366942805253810380628055...)i$.
I gather that Wolfram is instucted to choose one of the roots by some criteria that in this case leads to the exclusion of the real root. Maybe browsing the Wolfram site may help understanding what these criteria are. (My guess is that it outputs the root $\alpha=re^{i\theta}$ with smaller $\theta$ in the range $[0,2\pi)$.)
When non-computers calculate the cube root of (-8), we can think of it as $(-1*8)^{1/3}$ Then we have $-1*8^{1/3} = -1*2 = -2$
Wolfram is using the polar complex form of -8 = 8cis(π) Then the cube root of this is 2cis(π/3), which is 1 + i√3 (an alternate form on Wolfram)
Incidentally, if you take $(1 + i\sqrt3)^3$, you will get -8!
In the years since the question was asked and answered, Wolfram introduced the $\operatorname{Surd}(x,n)$ function (Mathematica 9 circa $2012$, then Alpha) to designate the real single-valued $n^{th}$ root of $x$.
For example $\sqrt[3]{-8}$ and $\sqrt[5]{-243}$ result directly in $-2$ and $-3$ respectively:
The $\operatorname{Cbrt}$ function discussed in Mark McClure's answer - which had changed behavior around the same time to return the real cube root by default - appears to be identical to $\operatorname{Surd}(\,\cdot\,,3)$:
Although the equation $x^3+8 = 0$ has actually three roots (real $-2$ and two conjugated complex roots) still the third root of $-8$ does not exist. Nth roots are defined only for nonnegative real values. Please, consider the $$(-8)^{1/3} = (-8)^{2/6}$$ that gives either $64^{1/6} = 2$ or $(\sqrt{ -8})^2$ = nonsense. Let $(-8)^{1/3} = (-8)^{2/6}$ you never get $(-2)$ as the result. Therefore Nth root exists only of nonnegative real numbers. | 2019-07-15T20:51:40 | {
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https://cathedralbaptist.com/7eujzh/qzuio6q.php?157bfb=adjacent-angles-which-are-not-in-a-linear-pair | # adjacent angles which are not in a linear pair
Consider the following figure in which a ray $$\overrightarrow{OP}$$ stand on the line segment $$\overline{AB}$$ as shown: The angles ∠POB and ∠POA are formed at O. ∠POB + … Answer/Explanation. Fill in the blanks: If two adjacent angles are supplementary, they form a _____. Linear Pair: Definition, Theorem & Example ... we can say that angle NSI and angle ISD also form a pair of adjacent angles. As the adjacent angles form a linear pair and they are supplementary. All adjacent angles are not linear, but all linear pairs are adjacent. Tell whether the angles are only adjacent, adjacent and form a linear pair, or not adjacent. True, if they are adjacent and share a vertex and one side. Angles ∠1 and ∠2 are non-adjacent angles. Which of the following CANNOT be true? Explanation for Linear Pair of Angles When the angle between the two lines is 180°, they form a straight angle. This statement is correct. But they are not a linear pair because they are not connected, and they do not share a common side. (1) Complementary angles that are not adjacent. They are supplementary because each angle is 90 degrees so they add up to 180 degrees. Related Questions to study. They share the same vertex and the same common side. D. The angles are congruent and are right angles. Two Adjacent Angle Can Be Complementary Too If They Add Up To 90°. Example: Find if following angles can make a linear pair. To identify whether the angles are adjacent or not, we must remember its basic properties that are … In the diagram below, ∠ABC and ∠DBE are supplementary since 30°+150°=180°, but they do not form a linear pair since they are not adjacent. O. Additionally, the ray OB is the common arm between these two angles. Those Adjacent Angles Are Complementary. Linear pair is a pair of adjacent angles where non-common side forms a straight line. Converse statement inverses statement Contrapositive statement conditional statement - edu … The angles are adjacent but not complementary. <3 and <4 form a linear pair. These Add Up To 180 Degrees (E And C Are Also Interiors). A Linear Pair Forms A Straight Angle Which Contains 180º, So You Have 2 Angles Whose Measures Add To 180, Which Means They Are Supplementary. A linear pair of angles is a pair of adjacent angles whose non common sides are opposite rays. Although they share a common side (PS) and a common vertex (S), they are not considered adjacent angles when they overlap like this. Answer: (b) Explanation : Definition of a linear pair of angles. Two angles form a linear pair. Yes. They have common vertex O. Last updated at Nov. 27, 2019 by Teachoo. In other words, if the non-common arms of a pair of adjacent angles are in a straight line, these angles make a linear pair. v. Angles which are neither complementary nor adjacent. If the two supplementary are adjacent to each other then they are called linear pair. As the vertically opposite angles are equal. angles (with their types) asked Aug 1, 2018 in Mathematics by vikashsoni (11.0k points) constructions; class-10; 0 votes. See our ‘Complementary and Supplementary Angles’ article for more details. One angle is obtuse and the other is acute. How to Find Adjacent Angles. Two angles are said to be linear if they are adjacent angles formed by two intersecting lines. Obviously, the larger angle ∠BAD is the sum of the two adjacent angles. At times, in geometry, the pair of angles are used. These linear pair of angles are always supplementary (both the angles sum up to 1800. The measure of a straight angle is 180 degrees, so a linear pair of angles must add up to 180 degrees. 265 views. In other words, if the non-common arms of a pair of adjacent angles are in a straight line, these angles make a linear pair. Let us take example of the angles, shown in following figure: ∠ COB and ∠ BOA have a common vertex, i.e. In this picture there is a flatscreen Toshiba television. Pair of adjacent whose measures add up … Linear Pair Of Angles : Two angles can be called as a linear pair, if they are adjacent angles formed by intersecting lines. Usually Found In Triangles And Other Polygons, Two Of The Sides That Meet At A Vertex Of The Polygon Are Called Adjacent Sides. What Are Adjacent Sides And Adjacent Angle? Sum of two adjacent supplementary = 180 o. They share a common vertex, but not a common side. Before you know all these pairs of angles there is another important concept which is called ‘angles on a straight line’. Definition: Two Lines That Meet At A Polygon Vertex. The measure of a straight angle is 180 degrees, so a linear pair of angles must add up to 180 degrees. Notice that angles OSN and ISD are not adjacent to each other. Pair of adjacent angles whose measures add up to form a straight angle is known as a linear pair. In a linear pair, the arms of the angles that are not always common or collinear i.e., they lie on a straight line. Vertically Opposite Angles: When two lines intersect, then the angles that are opposite one another at the intersection are called Vertically Opposite Angles. (c) -- Not sufficient. If Two Angles Form A Linear Pair, The Angles Are Supplementary. Solution: 110° + 70° = 180°Since the sum of these angles is equal to two right angles, so they can make a linear pair. Complementary angles that do not form a linear pair. Another way of defining them is: “two angles that share a side and a vertex, but do not share any interior points”. Angles in a linear pair which are not supplementary. The measure of rotation of a ray, when it is rotated about its endpoint is known as the angle formed by the ray between its initial and final position. They share a vertex and side, but donot overlap.A Linear Pair is twoadjacent angles whose non-common sides form opposite rays. In other words, they are angles that are side by side, or adjacent, sharing an "arm". Two angles are said to be linear if they are adjacent angles formed by two intersecting lines. iv. Two Obtuse Angles Can Be Adjacent Angles. Below are some pairs of complementary angles. Definition: Two angles that add up to 180°, Adjacent Angle – Definition, Examples & More. The two angles are said to be adjacent angles when they share the common vertex and side. Complementary angles are angle pairs that add up to exactly 90o. If two angles do not form a linear pair, then they are not supplementary. Example : Logical Equivalence Lateral Area . Common side. Adjacent Sides. Note: Two acute angles cannot make a linear pair because their sum will always be less than 180°. Solution: ∠ 1 and ∠ 2 are adjacent to each other∠ 2 and ∠ 3 are adjacent to each other, Solution: ∠ APD and ∠ DPC are adjacent to each other∠ DPC and ∠ CPB are adjacent to each other. A linear pair is two angles that add up to be 180o.A linear pair is two adjacent, supplementary angles.Adjacent means they share ONE ray.Supplementary means add … Adjacent Angles Share A Common Ray And Do Not Overlap. Linear Pair of Angles. Two adjacent angles forming a linear pair are in the ratio 7 : 5 find the angles. Two angles make a linear pair if their non-common arms are two opposite rays. Two angles are said to be supplementary if the sum of both the angles is 180 degrees. Two angles that overlap, one inside the other sharing a side and vertex in the figure on the right, the two angles ∠PSQ and ∠PSR overlap. iii. ii. In the figure above, the two angles ∠BAC and ∠CAD share a common side (the blue line segment AC). A pair of adjacent angles formed by intersecting lines is called a Linear Pair of Angles. Not necessarily true. Adjacent Angles Add Up To 180 Degrees. If Two Angles Form A Linear Pair, The Angles Are Supplementary. E. Both angles … The endpoints of the ray from the side of an angle are called the vertex of an angle. Note: Two acute angles cannot make a linear pair … Example: Find the pairs of adjacent angles in the following figure. A Linear Pair Forms A Straight Angle Which Contains 180º, So You Have 2 Angles Whose Measures Add To 180, Which Means They Are Supplementary. Bisect each of the two angles. (b) -- Not sufficient, because this definition does not require that the angles share a common ray. Can Two Adjacent Angles Be Supplementary? Adjacent angles can be a complementary angle or supplementary angle when they share the common vertex and side. (ii) Angles in a linear pair which are (iii) Complementary angles that do not form a not supplementary. Any two interior angles that share a common side are called the “adjacent interior angles” of the polygon or just “adjacent angles”. The angles in a linear pair are (a) complementary (b) supplementary (c) not adjacent angles (d) vertically opposite angles. Two angles together are angle pairs. Common vertex. Adjacent angles, Linear pair, Vertically opposite angles. A linear pair of angles has two defining characteristics: 1) the angles must be supplmentary; 2) The angles must be adjacent ; In the picture below, you can see two sets of angles. <1 and <2 <1 and <2 See the second picture. Any Two Angles That Add Up To 180 Degrees Are Known As Supplementary Angles. Start studying M54.2c. In the following picture, Р1 & Р2, Р2 & Р4, Р3 & Р4, and Р3 & Р4 are linear pairs. They share a common side, but not a common vertex. (i i i) When two lines intersect opposite angles are equal. If two angles are not adjacent, then they do not form a linear pair. Two adjacent angles are said to form a linear pair of angles, if their non-common arms are two opposite rays. The sum of angles on a straight line is 180 degree. The two angles will change so that they always add to 180°, In the figure above, the two angles ∠PQR and ∠JKL are supplementary because they always add to 180°. In A Right Triangle, The Altitude From A Right-angled Vertex Will Split The Right Angle Into Two Adjacent Angle; 30°+60°, 40°+50°, Etc. A And B Are Adjacent Angles. They are therefore termed ‘adjacent angles’. However, just because two angles are supplementary does not mean they form a linear pair. Vertical Angles Are Always Congruent, Which Means That They Are Equal. A linear pair of angles is formed when two lines intersect. Adjacent angles which are not in linear pair. The line through points A, B and C is a straight line. Adjacent Angles Are Angles That Come Out Of The Same Vertex. Supplementary angles a and b do not form linear pair. Two Adjacent Angle Can Be Supplementary Too, If They Add Upto 180°. Solution: 130° + 50° = 180°Since the sum of these angles is equal to two right angles, so they can make a linear pair. Both sets (top and bottom) are supplementary but only the top ones are linear pairs because these ones are also adjacent. What Are Adjacent Angles? Adjacent Angles Are Two Angles That Share A Common Vertex, A Common Side, And No Common Interior Points. i. Complementary angles that are not adjacent. | Definition & Examples - Tutors.com They have common side OB. In the figure, ∠1 and ∠3 are non-adjacent angles. ★★★ Correct answer to the question: 5. Hence, we can also say, that linear pair of angles is the adjacent angles whose non-common arms are basically opposite rays. If two angles form a linear pair, the angles are supplementary. What Are Adjacent Angles Or Adjacent Angles Definition? Example: If following angles make a linear pair, find the value of q. Smenevacuundacy and 26 more users found this answer helpful Here, these angles are in linear pair as. Draw the pairs of angles as described below. Adjacent Angles Are Two Angles That Share A Common Vertex, A Common Side, And No Common Interior Points. In Other Words, It Is Between A Right Angle And A Straight Angle. (i v) When two lines intersect opposite angles are supplementary. How to Square a Number in Java? Adjacent angles are angles that are side by side. Verify that the two bisecting rays are perpendicular to each other. They might not form a linear pair, like in a parallelogram. A. Enter your email address to subscribe to this blog and receive notifications of new posts by email. The angles in a linear pair are supplementary. In other words, when put together, the two angles form a right angle. So, In a linear pair, there are two angles who have. Two adjacent supplementary angles form a linear pair. Two angles make a linear pair if their non-common arms are two opposite rays. Non-common side makes a straight line or Sum of angles is 180°. Two angles in the same plane and with a common vertex and a common arm are called adjacent angles. (iv) Adjacent angles which are not in a linear pair (v) Angles which are neither complementary (vi) Angles in a linear pair which are nor adjacent. Since the non-adjacent sides of a linear pair form a line, a linear pair of angles is always supplementary. A polygon showing its interior angles, and a label pointing to two that are adjacent to another use of the term refers to the interior angles of polygons. Adjacent angles must be next to each other, not one on top of the other. Angles A and B are adjacent. Draw a pair of vertically opposite angles. Adjacent Angles, Linear Pair of angles, Vertically Opposite angles; Adjacent Angles. … They also share a common vertex (point A). Adjacent angles, often abbreviated as adj. One special relationship is called complementary angles. Sometimes, the measures of these angles form a special relationship. See the first picture below. Can Two Obtuse Angles Be Adjacent? 0 votes . If that is not possible, say why. asked Jul 15, 2019 in Class VI Maths by aditya23 (-2,145 points) Two adjacent angles forming a linear pair are in the ratio 7 : 5, find the angles. Now, we will learn more pairs of angles for grade 6 to grade 8 like linear, vertically opposite and adjacent angles here. ∠s, are angles that share a common vertex and edge but do not share any interior points. There are various kinds of pair of angles, like supplementary angles, complementary angles, adjacent angles, linear pairs of angles, opposite angles, etc. In the adjoining figure, ∠AOC and ∠BOC are two adjacent angles whose non-common arms OA and OB are two opposite rays, i.e., BOA is a line. Learn vocabulary, terms, and more with flashcards, games, and other study tools. A: Angles that are next to each other are known as adjacent angles, i.e., two angles with one common arm. Draw a linear pair of angles. Home » Geometry » Mathematics » Adjacent Angle – Definition, Examples & More. Solution: (7q – 46)° + (3q + 6)° = 180°Or, 10q – 40 = 180°Or, 10q = 180°+ 40 = 220°Or, q = 220° ÷ 10 = 22°. If two angles are not adjacent, then they do not form a linear pair. 5. : Math.pow() Method, Examples & More, Derivative Of Tangent – Slope, Derivative & More. vi. Here the word adjacent is used in its ordinary English meaning of “next to each other”. ∠POB and ∠POA are adjacent angles and they are supplementary i.e. When we say common vertex and a common side, we mean that the vertex point and the side are shared by the two angles. When two lines intersect, the angles in a linear pair have sides that are opposite rays; but the vertical angles formed by the two intersecting lines also have sides that are opposite rays. In this article, we are going to discuss the definition of adjacent angles and vertical angles in detail. C. The angles are supplementary, and the nonadjacent rays are opposite rays. … 1 answer. On the other hand, two right angles will always make a linear pair as their sum is equal to 180°. Adjacent angles have Common vertex Common side No overlap of angles Let’s look at some examples Adjacent angles are two angles that have a common vertex and a common side. Bisect each of the two angles. The pair of adjacent angles here are constructed on a line segment, but not all adjacent angles are linear. Linear pair of angles - definition Two angles are said to be in a linear pair if they are adjacent to each other, lie on the same side of the line and the sum of their measures is 1 8 0 o. Converse statement inverses statement… Get the answers you need, now! Learn what is linear pair of angles. It can also be said that angles of the linear pair are always supplementary to each other. Linear pair forms two supplementary angles. Hence, ∠ COB and ∠ BOA are the pairs of adjacent angles. This television has a pair of supplementary angles that are not a linear pair (one green, one blue). linear pair. Note: ∠ APD and ∠ CPB are not adjacent to each other, because they don’t have a common arm in spite of having a common vertex. The vertex of an angle is the endpoint of the rays that form the sides of the angle. B. | 2021-04-16T01:52:17 | {
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http://mathematica.stackexchange.com/questions/31470/how-to-plot-logarithmic-scales/31471 | # How to plot logarithmic scales
How can a simple logarithmic number line be drawn between any 2 integer values?
The closest function I found in the documentation is LogLinearPlot[] and I've been racking my brains trying to figure out how to do this with no luck...
-
## 4 Answers
One way is to plot the function 0 against a log axis.
LogLogPlot[0, {t, 1, 12}, Axes -> {True, False}, Ticks -> {Range[12]}]
or, changing the numbers
LogLogPlot[0, {t, 64, 96}, Axes -> {True, False}, Ticks -> {Range[64, 96]}]
The Axis function turns off the vertical axis (because you just want the number line) and the Ticks specifies where you want the tick marks. As a further example (to see the appropriate syntax), here is
LogLogPlot[0, {t, 1.07, 1.44}, Axes -> {True, False},
Ticks -> {{1.07, 1.15, 1.20, 1.29, 1.38, 1.44}}]
Note the double parentheses in the Ticks list. This is because Ticks is really a list of x-ticks and y-ticks (but since in this case, we aren't plotting any y's, so it is empty).
-
What if instead of a range of integers there is a list of irrational numbers? As LogLogPlot generates a plot of 0 as function of t from 1 to 1.5 for example, the number list doesn't seem to fit here. FindDivisions didn't help either so I guess it has to do with setting the ticks according to my defined list. I tried Ticks -> {1.07,1.15,1.20,1.29,1.38,1.44}, also with List but this doesn't work. – Bo C. Sep 1 '13 at 1:12
I added this example, above. Note the double parentheses in the Ticks list. This is because Ticks is really a list of x-ticks and y-ticks (but since in this case, we aren't plotting any y's, it's empty). – bill s Sep 1 '13 at 1:37
And another way:
logLine[min_, max_] := Module[{lines, labels},
lines = Line[{{Log[#], -1}, {Log[#], 1}}] & /@ Range[min, max];
labels = Text[#, {Log[#], 1.7}] & /@ Range[min, max];
Graphics[{
labels,
Line[{{Log[min], 0}, {Log[max], 0}}],
lines
}, AspectRatio -> 1/10
]
]
We have then that logLine[1,12] yields
To plot an arbitrary range we could use the following function:
logLineRange[range_] := Module[{lines, labels},
lines = Line[{{Log[#], -1}, {Log[#], 1}}] & /@ range;
labels = Text[#, {Log[#], 1.7}] & /@ range;
Graphics[{
labels,
Line[{{Log[Min[range]], 0}, {Log[Max[range]], 0}}],
lines
}, AspectRatio -> 1/10
]
]
Having defined that function, we can then do this:
logLineRange[{1.07, 1.15, 1.20, 1.29, 1.38, 1.44}]
-
Mathematica can't seem to plot this; tried replacing min and max with 1 and 12, both with and without the _ sign, only in the first line then all over the code - it returns nothing. I'm also interested in a version where the range of integers is replaced with a list of irrational numbers - for example 1.07,1.15,1.20,1.29,1.38,1.44. – Bo C. Sep 1 '13 at 8:59
@Bogdan The code block is just the definition for the function. To actually plot a scale you write logLine[1,12] for example. I added another function which takes an arbitrary range, you can have both those functions defined simultaneously as they take a different number of arguments. I added an example of how to use the second function as well. – C. E. Sep 1 '13 at 19:05
Thank you, everything is clear. How come your first example doesn't work for 1 to 10? Try logLine[1,10] – Bo C. Sep 2 '13 at 15:45
@Bogdan It works for me, don't know why it would not work for you. :/ – C. E. Sep 2 '13 at 15:49
Very strange, indeed. The Holy Restart solved it. As a side note, your second code can also be used for plotting a specific range with logLineRange[Range[1,12]] – Bo C. Sep 2 '13 at 17:26
Create unit-sized number lines with tick mapping function f for a list of values vals:
numberLine[f_, vals_List] :=
With[{pos = Rescale[f /@ vals], tick = 1/50},
Graphics[{Thick, Line[{{0, 0}, {1, 0}}],
MapThread[
{Line[{{#1, -tick}, {#1, tick}}],
Text[#2, {#1, 2 tick}]} &, {pos, vals}]},
PlotRange -> {{-2 tick, 1 + 2 tick}, {3 tick, -tick}}]]
GraphicsColumn[{numberLine[Log, Range[12]],
numberLine[Identity, Range[64, 96]]}]
EDIT: Apparently both number lines in the original question are logarithmic, while mine above are logarithmic and linear. This is easy to fix, of course.
-
Yes it is. On the last line, replace Identity with Log. I would also like to use your code to plot a version where the range of integers is replaced with a list of irrational numbers - for example 1.07,1.15,1.20,1.29,1.38,1.44. – Bo C. Sep 1 '13 at 9:04
@Bogdan You can replace vals with any list of numeric quantity; for instance, {1, 2, 3, 5, 7, 11, E, Pi, E^2, Pi^2, Sqrt[2], 5 Sin[1], 40/7}. Also, you can use Reals, such as 1.07. There might be some issues with this (solvable by Hold, or multipliers used with tick I believe), but for simpler expressions, it's just fine. – kirma Sep 1 '13 at 9:16
I can't seem to make this work; just replaced vals with the list - only in the first line then again in the second. I'm still searching the documentation to figure out what's wrong. This is the first time i use Mathematica; different versions of the code are my very best tutorial. Thank you. – Bo C. Sep 1 '13 at 10:13
@Bogdan Run Clear[numberLine], re-evaluate above definition of numberLine and try out numberLine[Identity, {1, 2, 3, 5, 7, 11, E, Pi, E^2, Pi^2, Sqrt[2], 5 Sin[1], 40/7}]. Maybe I was a bit vague. For beginners, Range[12], for instance, generates list corresponding to {1, 2, 3, ..., 12}. :) – kirma Sep 1 '13 at 11:02
Rescale[Log@v](b-a)+a will space any list v of positive values logarithmically over the interval [a,b].
v = Range[12]; Transpose@{N@Rescale[Log@v]*11+1, v}
{1., 1}
{4.06837, 2}
{5.86326, 3}
{7.13674, 4}
{8.12454, 5}
{8.93163, 6}
{9.61401, 7}
{10.2051, 8}
{10.7265, 9}
{11.1929, 10}
{11.6148, 11}
{12., 12}
v = Range[64,96,4]; Transpose@{N@Rescale[Log@v]*32+64, v}
{64., 64}
{68.7846, 68}
{73.2956, 72}
{77.5627, 76}
{81.6109, 80}
{85.4615, 84}
{89.1329, 88}
{92.6411, 92}
{96., 96}
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http://mathhelpforum.com/calculus/99864-vectors-application-3-a.html | # Math Help - Vectors Application 3
1. ## Vectors Application 3
Find the length of the median AM in the triangle ABC, for the points A(2,1.5,-4), B (3,-4,2) and C(1,3,-7), then find the distance from A to the centroid of the triangle.
Ok, so this is what I did so far, but I am stuck as to where to go to find the distance from A to the centroid of the triangle, plus my answer for the first part does not seem to be correct with the back.
BC Midpoint=(-2,7,-9)
----------
2
=(-1,7/2,-9/2)
AM=(-3,2,-0.5)
=sqrt (13.5)
that's my answer for part A, above, but the answer in the back says 2.5
For part b, I found the midpoint of each AB, BC, and AC.. but I am not sure how or where to go from here.
2. Hello, skeske1234!
(a) Find the length of the median $AM$ in $\Delta ABC$
with vertices: . $A(2,\tfrac{3}{2},-4),\;B (3,-4,2),\;C(1,3,-7)$
The midpoint is: . $\left(\frac{x_1 {\color{red}+}x_2}{2},\;\frac{y_1 {\color{red}+} y_2}{2},\;\frac{z_1{\color{red}+} z_2}{2}\right)$
$M$, the midpoint of $BC$ is: . $\left(\frac{3+1}{2},\;\frac{\text{-}4+3}{2},\;\frac{2-7}{2}\right) \;=\;\left(2,\;\text{-}\frac{1}{2},\;\text{-}\frac{5}{2}\right)$
Then: . $AM \;\;=\;\;\sqrt{(2-2)^2 + \left(\text{-}\frac{1}{2} - \frac{3}{2}\right)^2 + \left(\text{-}\frac{5}{2}+4\right)^2} \;\;=\;\;\sqrt{(0)^2 + (\text{-}2)^2 + \left(\frac{3}{2}\right)^2}$
. . . . . . . . $= \;\;\sqrt{0 + 5 + \frac{9}{4}} \;\;=\;\;\sqrt{\frac{25}{4}} \;\;=\;\;\frac{5}{2}$
(b) Find the distance from $A$ to the centroid of the triangle.
The centroid is $\tfrac{2}{3}$ of the way from $A\left(2,\:\tfrac{3}{2},\:-4\right)$ to $M\left(2,\:-\tfrac{1}{2},\:-\tfrac{5}{2}\right)$
Can you find it?
3. How do you know that the centroid is 2/3 of the way from A to midpoint? Is that the same for all centroids of a triangle in general? so from B to midpoint it would be 2/3 way for centroid, same with C?
4. Originally Posted by skeske1234
How do you know that the centroid is 2/3 of the way from A to midpoint? Is that the same for all centroids of a triangle in general? so from B to midpoint it would be 2/3 way for centroid, same with C?
That is a standard theorem: The medians of a triangle are concurrent and the distance from a vertex to that point is 2/3 of the length of the median
That theorem has an easy but messy proof using vectors. | 2015-02-01T08:25:23 | {
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https://maet.eat.kmutnb.ac.th/v0uuvx85/inner-product-of-complex-vectors-8f9b90 | inner product of complex vectors
Definition: The norm of the vector is a vector of unit length that points in the same direction as .. To motivate the concept of inner prod-uct, think of vectors in R2and R3as arrows with initial point at the origin. Example 3.2. two. Inner product spaces generalize Euclidean spaces (in which the inner product is the dot product, also known as the scalar product) to vector spaces of any (possibly infinite) dimension, and are studied in functional analysis. Of course if imaginary component is 0 then this reduces to dot product in real vector space. If the dimensions are the same, then the inner product is the trace of the outer product (trace only being properly defined for square matrices). There are many examples of Hilbert spaces, but we will only need for this book (complex length-vectors, and complex scalars). I see two major application of the inner product. �J�1��Ι�8�fH.UY�w��[�2��. A = [1+i 1-i -1+i -1-i]; B = [3-4i 6-2i 1+2i 4+3i]; dot (A,B) % => 1.0000 - 5.0000i A (1)*B (1)+A (2)*B (2)+A (3)*B (3)+A (4)*B (4) % => 7.0000 -17.0000i. Definition: The Inner or "Dot" Product of the vectors: , is defined as follows.. v|v = (v∗ x v∗ y v∗ z)⎛ ⎜⎝vx vy vz ⎞ ⎟⎠= |vx|2+∣∣vy∣∣2+|vz|2 (2.7.3) (2.7.3) v | v = ( v x ∗ v y ∗ v z ∗) ( v x v y v z) = | v x | 2 + | v y | 2 + | v z | 2. function y = inner(a,b); % This is a MatLab function to compute the inner product of % two vectors a and b. H�cf fc����ǀ |�@Q�%�� �C�y��(�2��|�x&&Hh�)��4:k������I�˪��. If both are vectors of the same length, it will return the inner product (as a matrix). Date . The Gelfand–Naimark–Segal construction is a particularly important example of the use of this technique. Nicholas Howe on 13 Apr 2012 Test set should include some column vectors. I want to get into dirac notation for quantum mechanics, but figured this might be a necessary video to make first. The outer product "a × b" of a vector can be multiplied only when "a vector" and "b vector" have three dimensions. Definition: The distance between two vectors is the length of their difference. �,������E.Y4��iAS�n�@��ߗ̊Ҝ����I���̇Cb��w��� Several problems with dot products, lengths, and distances of complex 3-dimensional vectors. 1 Real inner products Let v = (v 1;:::;v n) and w = (w 1;:::;w n) 2Rn. (Emphasis mine.) H�l��kA�g�IW��j�jm��(٦)�����6A,Mof��n��l�A(xГ� ^���-B���&b{+���Y�wy�{o������hC���w����{�|BQc�d����tw{�2O_�ߕ$߈ϦȦOjr�I�����V&��K.&��j��H��>29�y��Ȳ�WT�L/�3�l&�+�� �L�ɬ=��YESr�-�ﻓ�$����6���^i����/^����#t���! �X"�9>���H@ There are many examples of Hilbert spaces, but we will only need for this book (complex length vectors, and complex scalars). this special inner product (dot product) is called the Euclidean n-space, and the dot product is called the standard inner product on Rn. Then the following laws hold: Orthogonal vectors. A set of vectors in is orthogonal if it is so with respect to the standard complex scalar product, and orthonormal if in addition each vector has norm 1. 90 180 360 Go. Then their inner product is given by Laws governing inner products of complex n-vectors. Since vector_a and vector_b are complex, complex conjugate of either of the two complex vectors is used. Definition A Hermitian inner product on a complex vector space V is a function that, to each pair of vectors u and v in V, associates a complex number hu,vi and satisfies the following axioms, for all u, v, w in V and all scalars c: 1. hu,vi = hv,ui. For vectors with complex entries, using the given definition of the dot product would lead to quite different properties. And I see that this definition makes sense to calculate "length" so that it is not a negative number. Consider the complex vector space of complex function f (x) ∈ C with x ∈ [0,L]. How to take the dot product of complex vectors? ]��̷QD��3m^W��f�O' SVG AI EPS Show. There is no built-in function for the Hermitian inner product of complex vectors. x, y: numeric or complex matrices or vectors. This number is called the inner product of the two vectors. Definition: The length of a vector is the square root of the dot product of a vector with itself.. An inner product between two complex vectors, $\mathbf{c}_1 \in \mathbb{C}^n$ and $\mathbf{c}_2 \in \mathbb{C}^n$, is a bi-nary operation that takes two complex vectors as an input and give back a –possibly– complex scalar value. This generalization is important in differential geometry: a manifold whose tangent spaces have an inner product is a Riemannian manifold, while if this is related to nondegenerate conjugate symmetric form the manifold is a pseudo-Riemannian manifold. An inner product is a generalization of the dot product.In a vector space, it is a way to multiply vectors together, with the result of this multiplication being a scalar.. More precisely, for a real vector space, an inner product satisfies the following four properties. Show that the func- tion defined by is a complex inner product. for any vectors u;v 2R n, defines an inner product on Rn. this section we discuss inner product spaces, which are vector spaces with an inner product defined on them, which allow us to introduce the notion of length (or norm) of vectors and concepts such as orthogonality. Good, now it's time to define the inner product in the vector space over the complex numbers. Remark 9.1.2. �E8N߾+! I was reading in my textbook that the scalar product of two complex vectors is also complex (I assuming this is true in general, but not in every case). The term "inner product" is opposed to outer product, which is a slightly more general opposite. Inner Product. The inner productoftwosuchfunctions f and g isdefinedtobe f,g = 1 For each vector u 2 V, the norm (also called the length) of u is deflned as the number kuk:= p hu;ui: If kuk = 1, we call u a unit vector and u is said to be normalized. The inner product "ab" of a vector can be multiplied only if "a vector" and "b vector" have the same dimension. Algebraically, the dot product is the sum of the products of the corresponding entries of the two sequ We can call them inner product spaces. The first usage of the concept of a vector space with an inner product is due to Giuseppe Peano, in 1898. 3. EXAMPLE 7 A Complex Inner Product Space Let and be vectors in the complex space. An inner product on V is a map For complex vectors, the dot product involves a complex conjugate. In other words, the product of a by matrix (a row vector) and an matrix (a column vector) is a scalar. Inner products. We then define (a|b)≡ a ∗ ∗ 1b + a2b2. The length of a complex … Defining an inner product for a Banach space specializes it to a Hilbert space (or inner product space''). The inner product (or dot product'', or scalar product'') is an operation on two vectors which produces a scalar. For complex vectors, we cannot copy this definition directly. Minkowski space has four dimensions and indices 3 and 1 (assignment of "+" and "−" to them differs depending on conventions). As an example, consider this example with 2D arrays: Format. From two vectors it produces a single number. Real and complex inner products We discuss inner products on nite dimensional real and complex vector spaces. Another example is the representation of semi-definite kernels on arbitrary sets. Inner product of two arrays. Question: 4. Copy link. For complex vectors, the dot product involves a complex conjugate. For complex vectors, we cannot copy this definition directly. A vector space can have many different inner products (or none). The inner product and outer product should not be confused with the interior product and exterior product, which are instead operations on vector fields and differential forms, or more generally on the exterior algebra. Length of a complex n-vector. a complex inner product space $\mathbb{V}, \langle -,- \rangle$ is a complex vector space along with an inner product Norm and Distance for every complex inner product space you can define a norm/length which is a function Conjugate symmetry: $$\inner{u}{v}=\overline{\inner{v}{u}}$$ for all $$u,v\in V$$. For N dimensions it is a sum product over the last axis of a and the second-to-last of b: numpy.inner: Ordinary inner product of vectors for 1-D arrays (without complex conjugation), in higher dimensions a sum product over the last axes. Definition: The distance between two vectors is the length of their difference. a complex inner product space $\mathbb{V}, \langle -,- \rangle$ is a complex vector space along with an inner product Norm and Distance for every complex inner product space you can define a norm/length which is a function product. Solution We verify the four properties of a complex inner product as follows. Usage x %*% y Arguments. The test suite only has row vectors, but this makes it rather trivial. 1. Several problems with dot products, lengths, and distances of complex 3-dimensional vectors. However for the general definition (the inner product), each element of one of the vectors needs to be its complex conjugate. If a and b are nonscalar, their last dimensions must match. An inner product on is a function that associates to each ordered pair of vectors a complex number, denoted by , which has the following properties. Although we are mainly interested in complex vector spaces, we begin with the more familiar case of the usual inner product. Inner (or dot or scalar) product of two complex n-vectors. An inner product, also known as a dot product, is a mathematical scalar value representing the multiplication of two vectors. If the dot product of two vectors is 0, it means that the cosine of the angle between them is 0, and these vectors are mutually orthogonal. The dot product of two complex vectors is defined just like the dot product of real vectors. The reason is one of mathematical convention - for complex vectors (and matrices more generally) the analogue of the transpose is the conjugate-transpose. The Dot function does tensor index contraction without introducing any conjugation. A complex vector space with a complex inner product is called a complex inner product space or unitary space. Let and be two vectors whose elements are complex numbers. Defining an inner product for a Banach space specializes it to a Hilbert space (or inner product space''). The inner product "ab" of a vector can be multiplied only if "a vector" and "b vector" have the same dimension. complex-numbers inner-product-space matlab. All . So we have a vector space with an inner product is actually we call a Hilbert space. We can complexify all the stuff (resulting in SO(3, ℂ)-invariant vector calculus), although we will not obtain an inner product space. numpy.inner¶ numpy.inner (a, b) ¶ Inner product of two arrays. To verify that this is an inner product, one … 2. hu+v,wi = hu,wi+hv,wi and hu,v +wi = hu,vi+hu,wi. The Inner Product The inner product (or dot product'', or scalar product'') is an operation on two vectors which produces a scalar. When you see the case of vector inner product in real application, it is very important of the practical meaning of the vector inner product. Generalizations Complex vectors. Positivity: where means that is real (i.e., its complex part is zero) and positive. When a vector is promoted to a matrix, its names are not promoted to row or column names, unlike as.matrix. In the above example, the numpy dot function is used to find the dot product of two complex vectors. The Inner Product The inner product (or dot product'', or scalar product'') is an operation on two vectors which produces a scalar. (1.4) You should confirm the axioms are satisfied. Inner products on R defined in this way are called symmetric bilinear form. Inner product of two vectors. I don't know if there is a built in function for this, but you can implement your own: complexInner[a_, b_] := Conjugate[a].b This conjugates the first argument; you could in the same manner conjugate the second argument instead. In fact, every inner product on Rn is a symmetric bilinear form. H�lQoL[U���ކ�m�7cC^_L��J� %�D��j�7�PJYKe-�45$�0'֩8�e֩ٲ@Hfad�Tu7��dD�l_L�"&��w��}m����{���;���.a*t!��e�Ng���р�;�y���:Q�_�k��RG��u�>Vy�B�������Q��� ��P*w]T� L!�O>m�Sgiz���~��{y��r�����r�����K��T[hn�;J�]���R�Pb�xc ���2[��Tʖ��H���jdKss�|�?��=�ب(&;�}��H$������|H���C��?�.E���|0(����9��for� C��;�2N��Sr�|NΒS�C�9M>!�c�����]�t�e�a�?s�������8I�|OV�#�M���m���zϧ�+��If���y�i4P i����P3ÂK}VD{�8�����H��5�a��}0+�� l-�q[��5E��ت��O�������'9}!y��k��B�Vضf�1BO��^�cp�s�FL�ѓ����-lΒy��֖�Ewaܳ��8�Y���1��_���A��T+'ɹ�;��mo��鴰����m����2��.M���� ����p� )"�O,ۍ�. CC BY-SA 3.0. The dot product of two complex vectors is defined just like the dot product of real vectors. Each of the vector spaces Rn, Mm×n, Pn, and FI is an inner product space: 9.3 Example: Euclidean space We get an inner product on Rn by defining, for x,y∈ Rn, hx,yi = xT y. 3. . An inner product is a generalized version of the dot product that can be defined in any real or complex vector space, as long as it satisfies a few conditions.. The properties of inner products on complex vector spaces are a little different from thos on real vector spaces. 1. . The existence of an inner product is NOT an essential feature of a vector space. 2. . Definition Let be a vector space over .An inner product on is a function that associates to each ordered pair of vectors a complex number, denoted by , which has the following properties. In math terms, we denote this operation as: If we take |v | v to be a 3-vector with components vx, v x, vy, v y, vz v z as above, then the inner product of this vector with itself is called a braket. ;x��B�����w%����%�g�QH�:7�����1��~$y�y�a�P�=%E|��L|,��O�+��@���)��$Ϡ�0>��/C� EH �-��c�@�����A�?������ ����=,�gA�3�%��\�������o/����౼B��ALZ8X��p�7B�&&���Y�¸�*�@o�Zh� XW���m�hp�Vê@*�zo#T���|A�t��1�s��&3Q拪=}L��$˧ ���&��F��)��p3i4� �Т)|��q���nӊ7��Ob�$5�J��wkY�m�s�sJx6'��;!����� Ly��&���Lǔ�k'F�L�R �� -t��Z�m)���F�+0�+˺���Q#�N\��n-1O� e̟%6s���.fx�6Z�ɄE��L���@�I���֤�8��ԣT�&^?4ր+�k.��$*��P{nl�j�@W;Jb�d~���Ek��+\m�}������� ���1�����n������h�Q��GQ�*�j�����B��Y�m������m����A�⸢N#?0e�9ã+�5�)�۶�~#�6F�4�6I�Ww��(7��]�8��9q���z���k���s��X�n� �4��p�}��W8��v�v���G share. Share a link to this question. 164 CHAPTER 6 Inner Product Spaces 6.A Inner Products and Norms Inner Products x Hx , x L 1 2 The length of this vectorp xis x 1 2Cx 2 2. ⟩ factors through W. This construction is used in numerous contexts. A Hermitian inner product < u_, v_ > := u.A.Conjugate [v] where A is a Hermitian positive-definite matrix. The Inner Product The inner product (or dot product'', or scalar product'') is an operation on two vectors which produces a scalar. Ordinary inner product of vectors for 1-D arrays (without complex conjugation), in higher dimensions a sum product over the last axes. ^��t�Q��#��=o�m�����f���l�k�|�yR��E��~ �� �lT�8���6�c�|H� �%8Dxx&\aM�q{�Z�+��������6�$6�$�'�LY������wp�X20�f��w�9ׁX�1�,Y�� In Euclidean geometry, the dot product of the Cartesian coordinates of two vectors is widely used. More abstractly, the outer product is the bilinear map W × V∗ → Hom(V, W) sending a vector and a covector to a rank 1 linear transformation (simple tensor of type (1, 1)), while the inner product is the bilinear evaluation map V∗ × V → F given by evaluating a covector on a vector; the order of the domain vector spaces here reflects the covector/vector distinction. This ensures that the inner product of any vector … NumPy Linear Algebra Exercises, Practice and Solution: Write a NumPy program to compute the inner product of vectors for 1-D arrays (without complex conjugation) and in higher dimension. A row times a column is fundamental to all matrix multiplications. %PDF-1.2 %���� Details. Here the complex conjugate of vector_b is used i.e., (5 + 4j) and (5 _ 4j). This relation is commutative for real vectors, such that dot(u,v) equals dot(v,u). Defining an inner product for a Banach space specializes it to a Hilbert space (or inner product space''). Similarly, one has the complex analogue of a matrix being orthogonal. Hence, for real vector spaces, conjugate symmetry of an inner product becomes actual symmetry. Generalization of the dot product; used to defined Hilbert spaces, For the general mathematical concept, see, For the scalar product or dot product of coordinate vectors, see, Alternative definitions, notations and remarks. Examples and implementation. The Norm function does what we would expect in the complex case too, but using Abs, not Conjugate. Applied meaning of Vector Inner Product . Definition: The length of a vector is the square root of the dot product of a vector with itself.. Definition: The norm of the vector is a vector of unit length that points in the same direction as ..$\newcommand{\q}[2]{\langle #1 | #2 \rangle}$I know from linear algebra that the inner product of two vectors is 0 if the vectors are orthogonal. Product of vectors in Minkowski space is an example of indefinite inner product, although, technically speaking, it is not an inner product according to the standard definition above. Parameters a, b array_like. The inner product "ab" of a vector can be multiplied only if "a vector" and "b vector" have the same dimension. .$\begingroup$The meaning of triple product (x × y)⋅ z of Euclidean 3-vectors is the volume form (SL(3, ℝ) invariant), that gets an expression through dot product (O(3) invariant) and cross product (SO(3) invariant, a subgroup of SL(3, ℝ)). Kuifeng on 4 Apr 2012 If the x and y vectors could be row and column vectors, then bsxfun(@times, x, y) does a better job. This ensures that the inner product of any vector with itself is real and positive definite. An inner product space is a special type of vector space that has a mechanism for computing a version of "dot product" between vectors. Purely algebraic statements (ones that do not use positivity) usually only rely on the nondegeneracy (the injective homomorphism V → V∗) and thus hold more generally. Inner product spaces over the field of complex numbers are sometimes referred to as unitary spaces. In mathematics, the dot product or scalar product is an algebraic operation that takes two equal-length sequences of numbers, and returns a single number. The inner productoftwosuchfunctions f and g isdefinedtobe f,g … If the dot product is equal to zero, then u and v are perpendicular. Suppose We Have Some Complex Vector Space In Which An Inner Product Is Defined. [/������]X�SG�֍�v^uH��K|�ʠDŽ�B�5��{ҸP��z:����KW�h���T>%�\���XX�+�@#�Ʊbh�m���[�?cJi�p�؍4���5~���4c�{V��*]����0Bb��܆DS[�A�}@����x=��M�S�9����S_�x}�W�Ȍz�Uή����Î���&�-*�7�rQ����>�,$�M�x=)d+����U���� ��հ endstream endobj 70 0 obj << /Type /Font /Subtype /Type1 /Name /F34 /Encoding /MacRomanEncoding /BaseFont /Times-Bold >> endobj 71 0 obj << /Filter /FlateDecode /Length 540 /Subtype /Type1C >> stream The inner product is more correctly called a scalar product in this context, as the nondegenerate quadratic form in question need not be positive definite (need not be an inner product). (����L�VÖ�|~���R��R�����p!۷�Hh���)�j�(�Y��d��ݗo�� L#��>��m�,�Cv�BF��� �.������!�ʶ9��\�TM0W�&��MY�>�i�엑��ҙU%0���Q�\��v P%9�k���[�-ɛ�/�!\�ے;��g�{иh�}�����q�:!NVز�t�u�hw������l~{�[��A�b��s���S�l�8�)W1���+D6mu�9�R�g،. Returns out ndarray. An interesting property of a complex (hermitian) inner product is that it does not depend on the absolute phases of the complex vectors. Ordinary inner product of vectors for 1-D arrays (without complex conjugation), in higher dimensions a sum product over the last axes. And so this needs a little qualifier. Let a = and a1 b = be two vectors in a complex dimensional vector space of dimension . We can complexify all the stuff (resulting in SO(3, ℂ)-invariant vector calculus), although we will not obtain an inner product space. 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An inner product space is a special type of vector space that has a mechanism for computing a version of "dot product" between vectors. Note that the outer product is defined for different dimensions, while the inner product requires the same dimension. Let X, Y and Z be complex n-vectors and c be a complex number. In other words, the inner product or the vectors x and y is the product of the magnitude s of the vectors times the cosine of the non-reflexive (<=180 degrees) angle between them. A bar over an expression denotes complex conjugation; e.g., This is because condition (1) and positive-definiteness implies that, "5.1 Definitions and basic properties of inner product spaces and Hilbert spaces", "Inner Product Space | Brilliant Math & Science Wiki", "Appendix B: Probability theory and functional spaces", "Ptolemy's Inequality and the Chordal Metric", spectral theory of ordinary differential equations, https://en.wikipedia.org/w/index.php?title=Inner_product_space&oldid=1001654307, Short description is different from Wikidata, Articles with unsourced statements from October 2017, Creative Commons Attribution-ShareAlike License, Recall that the dimension of an inner product space is the, Conditions (1) and (2) are the defining properties of a, Conditions (1), (2), and (4) are the defining properties of a, This page was last edited on 20 January 2021, at 17:45. I was reading in my textbook that the scalar product of two complex vectors is also complex (I assuming this is true in general, but not in every case). Let X, Y and Z be complex n-vectors and c be a complex number. 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If a and b are nonscalar, their last dimensions must match dimensions must match R in! Useful alternative notation for inner products we discuss inner products on R defined in way.: numeric or complex n-tuple s, the vectors:, is defined as follows part is zero and! Vectors, we can not copy this definition directly 1b + a2b2 more familiar case of the needs. Without complex conjugation ), in higher dimensions a sum product over a complex inner in! Product of two complex vectors x, Y: numeric or complex n-tuple s, the definition changed! Matrix being orthogonal vector_b is used in numerous contexts v 2R n, defines an product... Dot ( v, u ) root of the Cartesian coordinates of two complex,! N-Vectors and c be a necessary video to make first section v is a complex space! 2. hu+v, wi = hu, v ) equals its complex conjugate ⟩ factors through W. this construction a... It will return the inner product space '' ) important example of the use of this technique,. U and v are perpendicular inner product is equal to zero, then: second vector provide the of! + 4j ) and positive of Hilbert spaces, conjugate symmetry of an inner is! Nonzero vector space with a complex vector space of two complex vectors is defined for different dimensions, while inner. The two complex vectors, the dot product of x and Y is square... Tensor index contraction without introducing any conjugation invented a useful alternative notation for quantum mechanics, but Abs... Wi and hu, wi+hv, wi g isdefinedtobe f, g = inner. Is defined as follows positive-definite matrix ( as a matrix ) at the origin first of... Two complex vectors, we can not copy this definition directly space, then: product the..., and be vectors and be a scalar, then this reduces to dot product of and... This definition makes sense to calculate length '' so that it is not an essential feature a... Will return the inner product for a Banach space specializes it to a Hilbert space ( ! Real vector space, then this reduces to dot product of the dot product of in! 3-Dimensional vectors n, defines an inner product of the vector space in which an inner product ) =,. Is straight commutative for real vector spaces ( without complex conjugation ), each element of one the! Copy this definition directly inner ( or inner is horizontal times vertical and shrinks,. To as unitary spaces norm function does what we would expect in the same length, it will return inner... Is widely used to quite different properties suppose we have some complex vector space with an inner product in vector... Why the inner product over the field of complex 3-dimensional vectors space )! Here the complex numbers actual symmetry vector is a vector space can have many different inner products each... The vectors u ; v 2R n, defines an inner product times and... Axioms are satisfied as follows space or unitary space as the length of a vector space defined in section. Matrix, its complex conjugate f ( x ) ∈ c with x ∈ [ 0, ]! Standard dot product is given by Laws governing inner products allow the introduction! Built-In function for the Hermitian inner product requires the same direction as now, we denote this operation as Generalizations..., each element of one of the usual inner product ), each element of one of the product. Course if imaginary component is 0 then this is a Hermitian inner product a finite dimensional vector involves... As an inner product for complex vector spaces, conjugate symmetry of inner. Products on R defined in this way are called symmetric bilinear form 2012 test set should some! Dot products, lengths, and complex vector space in which an inner.... Defined for different dimensions, while the inner or inner product a. And g isdefinedtobe f, g = 1 inner product is defined the! Bilinear form suppose we have a vector space involves the conjugate of the vector is a vector with..... _ 4j ) and positive examples of Hilbert spaces, we can not copy this definition directly semi-definite on. Than the conventional mathematical notation we have a vector space four properties of a vector space with inner..., lengths, and be vectors and be two vectors products of complex function f ( x ∈... '' is opposed to outer product is actually we call a Hilbert space ( . In higher dimensions a sum product over the field of complex numbers a scalar,:. Representation of semi-definite kernels on arbitrary sets component of the second vector the inner... Space over the last axes { R } \ ) equals dot (,. Y and Z be complex n-vectors this way are called symmetric bilinear form is 0 then this is straight for! Referred to as unitary spaces u, v +wi = hu, vi+hu,.... Provide the means of defining orthogonality between vectors ( zero inner product two. Definition is changed slightly vector spaces v_ >: = u.A.Conjugate [ v ] a! Operation as: Generalizations complex vectors this definition makes sense to calculate length '' so it. Dot products, lengths, and distances of complex vectors is used i.e., its complex part is zero and. The inner product of the concept of a vector with itself is real ( i.e., its names not! Vertical and shrinks down, outer is vertical times horizontal and expands out.! Y: numeric or complex matrices or vectors of a vector with itself this construction is used definition changed... The test suite only has row vectors, such that dot ( v, u ) definition: the function. Any vectors u ; v 2R n, defines an inner product space '' ) vector or angle. But we will only need for this book ( complex length-vectors, and complex scalars ), one the! | 2021-06-16T23:42:46 | {
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https://math.stackexchange.com/questions/2753982/how-many-different-ways-are-there-to-choose-5-items-from-12-distinct-items-if | # How many different ways are there to choose 5 items from 12 distinct items if…
How many different ways are there to choose 5 items from 12 distinct items if…
a. (5 pt.) the order of the items matters and repetition of items is not allowed?
b. (5 pt.) the order of the items matters and repetition of items is allowed?
c. (5 pt.) the order of the items does not matter and repetition of items is not allowed?
d. (5 pt.) the order of the items does not matter and repetition of items is allowed?
My Work
a) $P(12,5)$
b)$12^5$
c)$\binom{12}{5}$
Can you please verify my work
• I don't understand your reasoning for $b$. Hint: there are $12$ possibilities for the first choice, $12$ possibilities for the second choice, ... – lulu Apr 25 '18 at 23:39
• So, the b is $12^5$ right – tien lee Apr 25 '18 at 23:41
• Yes, that is right. For $d$, I think Stars and Bars is the way to go. – lulu Apr 25 '18 at 23:42
• Then, what about C – tien lee Apr 25 '18 at 23:45
• Your answers for $a,c$ are correct,. – lulu Apr 25 '18 at 23:46
a) Correct
b) No. You have 12 choices on the first item. 12 on the second. 12 on the third, fourth and fifth. So it should be $12^5$.
c) Correct
d) This is a bit tricky, I shall explain this below.
We have $C$ for combination, $P$ for permutation, and the first time I was taught this one (order does not matter + repetition allowed), my teacher called it $H$ although I am not sure if it is standard notation.
Anyway, the formula in general is $H_r^{n}=C_r^{n+r-1}$.
The reasoning behind this is as follows (for simplicity we shall take $H_5^{12}$ as in this case as an example):
Let the $12$ possible choices be $a,b,c,\dots,l$.
We don't care about the order here, we only care about how many of each item you choose, and with each choice, we can assign with it a unique "binary code".
For example, if you choose $aabbc$, then we write $0010010111111111$ where the first two $0$'s represent the two $a$'s, after the $1$ we have two $0$'s for the two $b$'s, and so on.
As another example, $abjkl$ is represented by $0101111111101010$.
In this way, each possible choice can be represented by a string of five $0$'s and eleven $1$'s for a total of $16$ numbers. So to count all possible choices, we simply have to choose the five places to insert the $0$'s.
This is given by $C_5^{16}$, hence $H_5^{12}=C_5^{16}$
• Isn't $C(16,4)$ – tien lee Apr 25 '18 at 23:57
• Why should it be $C_4^{16}$? – glowstonetrees Apr 26 '18 at 0:07
• @tienlee If the number of items of type $k$ is $x_k$, then we want the number of solutions of the equation $$x_1 + x_2 + x_3 + \cdots + x_{12} = 5$$ in the nonnegative integers. The number of such solutions is $$\binom{12 + 5 - 1}{12 - 1} = \binom{12 + 5 - 1}{5}$$ where the expression $\binom{12 + 5 - 1}{12 - 1}$ counts the number of ways of choosing which $12 - 1 = 11$ positions of $12 + 5 - 1$ positions required for five ones and eleven addition signs will be filled with addition signs and the expression $\binom{12 + 5 - 1}{5}$ comes from choosing the positions of the ones in the same row. – N. F. Taussig Apr 26 '18 at 0:19
I will leave you with a quick guide on how to tackle these problems.
Suppose you have $n$ distinct objects, and you want to select $r$ out of these $n$ objects. Then, if:
• Order is relevant and repetition is not allowed: look at the possible permutations of size $r$ which is given by $$P(n,r)=\frac{n!}{(n-r)!}$$ with $0\leq r \leq n.$
• Order is relevant and repetition is allowed: look at the possible arrangements of size $r$, which is given by $n^r$, with $n,r \geq 0.$
• Order is not relevant and repetitions are not allowed: look at the possible combinations of size $r,$ which is given by $$C(n,r)= \binom{n}{r}=\frac{n!}{(n-r)! r!}$$ with $0 \leq r \leq n.$
• Order is not relevant and repetition is allowed: look at the possible combinations with repetitions, which is given by $$\binom{n+r-1}{r}\ \$$ with $n,r \geq 0.$ | 2020-10-31T14:14:16 | {
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https://www.seamplex.com/wasora/realbook/real-010-2dfunctions.html | # Two-dimensional functions
• Difficulty: 010/100
• Author: jeremy theler
• Keywords: FUNCTION, :=, FILE_PATH, DATA, INTERPOLATION, nearest, rectangle, PRINT_FUNCTION, MIN, MAX, STEP, NUMBER, FILE, OUTPUT_FILE,
These examples illustrate the facilities wasora provides to interpolate two-dimensional functions. As shown in section ref{007-functions}, multidimensional functions may be defined using algebraic expressions, inline data, external files, vectors or dynamically-loaded routines. The main focus of this section is to illustrate the difference between algebraic and pointwise-defined two-dimensional functions, and to further illustrate the difference between a nearest-neighbor interpolation (based on a $$k$$-dimensional tree structure) and a rectangle interpolation (based on finite-elements-like shape functions for quadrangles).
## paraboloid.was
This input defines an algebraic function $$f(x,y)$$ and uses PRINT_FUNCTION to dump its contents (as three columns containing $$x$$, $$y$$ and $$f(x,y)$$ as shown in the terminal mimic) to the standard output. The range is mandatory because $$f(x,y)$$ is defined by an algebraic expression.
As far as the function $$f(x,y)$$ is concerned, $$a$$ and $$b$$ are taken as parameters. Even though they can change over time, the value they take is the value they have when $$f(x,y)$$ is evaluated. In this case, $$f(x,y)$$ is evaluated when the instruction PRINT_FUNCTION is executed, so the output is written with $$b=2$$.
a = 1
b = 1
f(x,y) := (x/a)^2 + (y/b)^2
b = 2
# range is mandatory as f(x,y) is algebraically-defined
PRINT_FUNCTION f MIN -b -b MAX b b STEP b/20 b/20
$wasora paraboloid.was > paraboloid.dat$ head paraboloid.dat
-2.000000e+00 -2.000000e+00 5.000000e+00
-2.000000e+00 -1.900000e+00 4.902500e+00
-2.000000e+00 -1.800000e+00 4.810000e+00
-2.000000e+00 -1.700000e+00 4.722500e+00
-2.000000e+00 -1.600000e+00 4.640000e+00
-2.000000e+00 -1.500000e+00 4.562500e+00
-2.000000e+00 -1.400000e+00 4.490000e+00
-2.000000e+00 -1.300000e+00 4.422500e+00
-2.000000e+00 -1.200000e+00 4.360000e+00
-2.000000e+00 -1.100000e+00 4.302500e+00
$gnuplot paraboloid.gp$
## nearest.was
This input defines a two-dimensional pointwise-defined function $$g(x,y)$$ using data inlined in the input file using the DATA keyword. As can be seen, expressions are allowed. However, they are evaluated at parse-time so references to variables should be avoided, as they will default to zero. For that end, numbers defined with the NUMBER keyword should be used. By default, pointwise-defined multidimensional functions are interpolated by the nearest neighbor algorithm (i.e. default is INTERPOLATION nearest). When calling PRINT_FUNCTION with no range, the definition points are printed. When a range is given, the function gets evaluated at the grid points.
In this case, first the function $$g(x,y)$$ is dumped into a file called g_def.dat with no range. Then, the function is dumped into a file called g_int.dat using a range and thus resulting in an interpolated output using nearest neighbors.
FUNCTION g(x,y) DATA {
0 0 1-1
0 1 1-0.5
0 2 1
1 0 1
1 1 1+0.25
1 2 1
2 0 1-0.25
2 1 1+0.25
2 2 1+0.5
}
# print g(x,y) at the definition points
PRINT_FUNCTION g FILE_PATH g_def.dat
# print g(x,y) at the selected range
# by defaults wasora interpolates using nearest neighbors
PRINT_FUNCTION g FILE_PATH g_int.dat MIN 0 0 MAX 2 2 STEP 0.05 0.05
$wasora nearest.was$ gnuplot nearest.gp
$ ## rectangle.was This time, a function $$h(x,y)$$ is defined by reading point-wise data from a file. This file is g_def.dat which are the definition points of the function $$g(x,y)$$ of the previous example. Note that when reading function data from a file, no expressions are allowed. Function $$h(x,y)$$ is interpolated using the rectangle method. The interpolated data is written in a file called h_int.dat, in which the function $$h(x,y)$$ is evaluated at the very same points the function $$g(x,y)$$ of the previous example was. # h(x,y) is equal to g(x,y) at the definition points, # but it is interpolated differently FUNCTION h(x,y) INTERPOLATION rectangle FILE_PATH g_def.dat # print h(x,y) at the selected range PRINT_FUNCTION h FILE_PATH h_int.dat MIN 0 0 MAX 2 2 STEP 0.05 0.05 $ cat g_def.dat
0.000000e+00 0.000000e+00 0.000000e+00
0.000000e+00 1.000000e+00 5.000000e-01
0.000000e+00 2.000000e+00 1.000000e+00
1.000000e+00 0.000000e+00 1.000000e+00
1.000000e+00 1.000000e+00 1.250000e+00
1.000000e+00 2.000000e+00 1.000000e+00
2.000000e+00 0.000000e+00 7.500000e-01
2.000000e+00 1.000000e+00 1.250000e+00
2.000000e+00 2.000000e+00 1.500000e+00
$wasora rectangle.was$ gnuplot rectangle.gp
$ ## scattered.was The last example of two-dimensional interpolation involves a pointwise-defined function~$$s(x,y)$$ using scattered data, i.e. not necessarily over a rectangular grid. In this case—and with no information about any underlying finite-element-like mesh—wasora can use either a nearest-neighbor interpolation or a Shepard-like inverse distance weighting. For the original Shepard method, the only parameter than can be tweaked is the exponent~$$p$$ of the distance in the weight~$$w_i=1/d_i^p$$. For the modified Shepard algorithm, the radius~$$r$$ of the nearest neighbors taken into account is to be provided. These nearest neighbors are found using a $$k$$-dimensional tree, that is a very efficient way of doing this task. For complex functions, all the alternatives should be investigated taking into account accuracy and code speed. # scattered multidimensional data may be interpolated # using a nearest-neighbor approach FUNCTION n(x,y) INTERPOLATION nearest DATA { 0 0 0 1 0 1 0 1 2 -0.5 0.5 3 -1 -1 2 0.75 0 1.5 0.25 0.25 1 } # another way of giving the same set of data VECTOR datax SIZE 7 DATA 0 1 0 -0.5 -1 0.75 0.25 VECTOR datay SIZE 7 DATA 0 0 1 0.5 -1 0 0.25 VECTOR dataz SIZE 7 DATA 0 1 2 3 2 1.5 1 # using shepard's interpolation method FUNCTION s(x,y) VECTORS datax datay dataz INTERPOLATION shepard SHEPARD_EXPONENT 4 # or using shepard's modified algorithm FUNCTION m(x,y) VECTORS datax datay dataz INTERPOLATION modified_shepard SHEPARD_RADIUS 2 # print the definition points PRINT_FUNCTION n FILE_PATH n_def.dat # print the different functions at the selected range PRINT_FUNCTION n s m FILE_PATH n_int.dat MIN -1 -1 MAX 1.5 1.5 STEP 0.05 0.05 $ wasora scattered.was
$gnuplot scattered.gp$
The figures illustrate how the multidimensional data interpolation scheme work for pointwise defined functions over scattered data. Nearest neighbors give constant values for each voronoi triangle whilst Shepard-based algorithms provide continuous surfaces.
## compwater.was
This example shows an extension of the example about saturated water in section ref{007-functions} by giving properties of compressed water as a function of pressure $$p$$ and temperature $$T$$. The file compwater.txt contains some properties of water as a function of temperature and pressure in a rectangular grid over the pressure-temperature space. The enthalpy is not contained in the file, but it can be computed from pressure $$p$$, the internal energy $$u(p,T)$$ and the specific volume $$v(p,T)$$ as
!bt [ h(p,t) = u(p,T) + p v(p,T) ] !et
FUNCTION v(p,T) FILE_PATH compwater.txt COLUMNS 2 1 4 INTERPOLATION rectangular
FUNCTION u(p,T) FILE_PATH compwater.txt COLUMNS 2 1 5 INTERPOLATION rectangular
FUNCTION h(p,T) = u(p,T) + p*v(p,T)
PRINT_FUNCTION h MIN 1e5 300 MAX 200e5 1000 STEP 2e5 4
$wasora compwater.was > compwater.dat$ gnuplot compwater.gp
\$
The figure shows the enthalpy of compressed water as a continuous function of $$p$$ and $$T$$ and the discrete experimental data. | 2018-12-12T19:24:29 | {
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https://physics.stackexchange.com/questions/500339/why-does-newtons-first-law-create-two-different-answers-to-this-question/500367 | # Why does Newton's first law create two different answers to this question?
I have been having some difficulty with a recent question. Take the following pulley system:
The three blocks have masses $$M$$, $$m_1$$, and $$m_2$$. All are subjected to a gravitational force $$g$$. The pulley and string are of negligible mass, and all surfaces are frictionless. The problem is (to paraphrase):
What magnitude of force $$F$$ is necessary for $$m_1$$ and $$m_2$$ to be motionless relative to $$M$$?
When I first solved this, I just considered Newton's first law for each of the three masses and added the additional conditions $$a_{xm_1}=a_{xM}$$, $$a_{xm_2}=a_{xM}$$, $$a_{ym_1}=0$$. After eliminating most of the variables, I ended up with $$F=g\frac{m_1}{m_2}(M+m_1)$$.
However, the textbook gives the answer $$F=g\frac{m_1}{m_2}(M+m_1+m_2)$$. In attempting to find the discrepancy, I solved the problem again somewhat differently: Let $$a=a_{xM}=a_{xm_1}=a_{xm_2}$$. Since $$m_1$$ does not accelerate vertically, $$T=m_1g$$. Since $$m_2a_{xm_2}=T=m_1g$$, we have that $$a=g\frac{m_1}{m_2}$$. Finally, since $$F$$ is pushing on $$M$$, which in turn is pushing on $$m_1$$ via its normal-force interaction, we have $$F=(M+m_1)a=g\frac{m_1}{m_2}(M+m_1)$$, the same answer I previously came to.
I asked my teacher about this problem to determine where my error occurred, and his reasoning was as follows: As with the earlier reasoning, $$a=g\frac{m_1}{m_2}$$. Since $$M$$, $$m_1$$, and $$m_2$$ are motionless relative to each other, we can treat the three as a single system and ignore internal forces:
Now, we simply have $$F=M_sa=g\frac{m_1}{m_2}(M+m_1+m_2)$$.
Both lines of reasoning are compelling, so my overall question is this: Which of these two answers is correct, and how is the other answer incorrect?
• It seems that I indeed forgot to account for the force on $M$ from the pulley. Adding that to the equations produces the correct result. I've accepted this answer, even though both are equally valid, since it is the earlier one. – LegionMammal978 Sep 6 '19 at 20:28
You have made a mistake in the equation $$F=(M+m_1)a$$ You have missed the fact that the tension is also works on M.! The sum of the two tensions add up and produce a force on the pulley. However, as the pulley is(should be) massless, the block M experiences a opposite force, which has a horizontal component of $$T√2*cos45° = T$$ So the equation should be $$F-T=(M+m_1)a$$ which then gives same answer as your teacher. | 2021-07-29T12:16:39 | {
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https://math.stackexchange.com/questions/3141354/show-a-specially-defined-matrix-is-positive-definite?noredirect=1 | # Show a specially defined matrix is positive definite
Let $$E_1, ..., E_n$$ be non empty finite sets. Show that the matrix $$A = (A_{ij})_{1 \leq i, j \leq n}$$ defined by $$A_{ij} = \dfrac{|E_i \cap E_j|}{|E_i \cup E_j|}$$, is positive semi-definite.
This is part 5 of Exercise 1 in http://members.cbio.mines-paristech.fr/~jvert/svn/kernelcourse/homework/2019mva/hw.pdf .
• If any two of the sets are empty then the corresponding coefficient is not defined. If the one set is disjoint from all the others, then its corresponding row and column are zero, so the matrix is not positive definite. – Servaes Mar 9 '19 at 18:53
• You are correct, I've corrected the problem accordingly. – SiXUlm Mar 9 '19 at 21:14
• Could the semi-definite positiveness be connected to the submodularity mentionned in the following answer : math.stackexchange.com/q/182384 ? – Jean Marie Mar 9 '19 at 23:09
• Let $E$ be the union of the $E_i$. For each $e \in E$, define the matrix $B_e = \left( \dfrac{\left[e \in E_i \cap E_j\right]}{\left|E_i \cup E_j\right|} \right)_{i, j \in \left[n\right]}$, where square brackets stand both for the Iverson bracket and for $\left[n\right] := \left\{1,2,\ldots,n\right\}$. Clearly, $A = \sum_{e \in E} B_e$. Are the $B_e$ still positive semidefinite? This would likely be easier if true. After all, for each $e \in E$, the matrix $B_e$ is just a blown-up matrix with entries $\dfrac{1}{\left|E_i \cup E_j\right|}$. – darij grinberg Mar 10 '19 at 3:01
• Even stronger: For any $e \in E$ and $t \in \mathbb{R}_+$, let $C_{e,t}$ be the matrix $\left(\left[e \in E_i \cap E_j\right] t^{\left|E_i \cup E_j\right|}\right)_{i, j \in \left[n\right]}$. Is $C_{e,t}$ nonnegative semidefinite for $0 < t < 1$ ? If it is, then so is $B_e$, since $B_e = \int_0^1 \left(C_{e,t}/t\right) dt$. – darij grinberg Mar 10 '19 at 3:08
Assume the sets $$E_i$$ are all non-empty. By deleting rows and columns of $$A$$ if needed one can reduce to the case where the $$E_i$$ are all distinct.
Write $$A_{ij}= \frac{|E_i\cap E_j|}{|E_i\cup E_j|} = \frac{|E_i\cap E_j|}{|E_i|+|E_j|-|E_i\cap E_j|} = \frac{e_{ij}}{e_i+e_j-e_{ij}}$$ where $$e_i$$ is shorthand for $$|E_i|$$ and so on. Then $$A_{ij}= \frac{e_{ij}}{e_i+e_j} \left( 1 + r_{ij} + r_{ij}^2 + r_{ij}^3 + \cdots\right)\tag{*}$$ where $$r_{ij}= e_{ij}/(e_i+e_j).$$ The distinctness hypothesis implies $$|r_{ij}|<1$$ and so the convergence of the geometric series in (*).
The matrix with entries $$e_{ij}$$ is a Gram matrix and hence positive semidefinite. The matrix with entries $$1/(e_i+e_j) = \int_0^\infty e^ {-x|E_i|} e^{-x|E_j|} \,dx$$ is also a Gram matrix so it too is psd. Schur's product theorem tells us that their elementwise product $$r_{ij}$$ is thus also psd. By Schur again, all the summand matrices in (*) are psd. So the sum, $$A$$, is psd.
Examples (where all the $$E_i$$ are equal, so all $$A_{ij}=1$$, say) show that positive semi definite cannot be replaced by positive definite as in an earlier form of the problem statement. I suppose that strict positive definiteness holds in the reduced case, but I don't see an argument. Numerical experiments show that the matrix $$(r_{ij})$$ is not strictly positive definite when (for example) the $$E_i$$ are all the 2-element subsets of $$[4]$$. The same experiments however do show $$A$$ strictly positive definite.
• Thanks for your solution. I've edited the question because I didn't state it correctly. – SiXUlm Mar 9 '19 at 21:28
• Thanks for the clarification. I've just found out that you talked about Schur product theorem. – SiXUlm Mar 9 '19 at 21:35
• [+1] Very smart ! – Jean Marie Mar 9 '19 at 22:37
• Why is $(e_{ij})=(|E_i\cap E_j|)$ a Gram matrix? – William McGonagall Feb 14 at 12:32
• @WilliamMcGonagall Because $e_{ij}=\langle v_i,v_j\rangle$, where $v_r$ denotes the vector whose $k$-th component is $1$ exactly when $k\in E_r$, and zero otherwise, and so on. – kimchi lover Feb 14 at 14:21
The following proof is inspired by my answer to math.stackexchange question https://math.stackexchange.com/q/1340405/, which in turn is inspired by the probabilistic method from extremal combinatorics. It is fully elementary and almost combinatorial ("almost" because it involves a simple limit argument at one point).
We begin with notations:
Let $$\mathbb{N}$$ be the set $$\left\{ 0,1,2,\ldots\right\}$$.
For each $$n\in\mathbb{N}$$, let $$\left[ n\right]$$ denote the set $$\left\{ 1,2,\ldots,n\right\}$$.
We also recall the following simple fact (Lemma 5 in my answer to math.stackexchange question https://math.stackexchange.com/q/1340405/):
Lemma 1. Let $$Q$$ be a finite totally ordered set. Let $$J$$ be a subset of $$Q$$. Let $$r\in J$$. Let $$S$$ be the set of all permutations of $$Q$$. Then, $$$$\left\vert \left\{ \sigma\in S\ \mid\ \sigma\left( r\right) >\sigma\left( j\right) \text{ for all }j\in J\setminus\left\{ r\right\} \right\} \right\vert =\dfrac{\left\vert S\right\vert }{\left\vert J\right\vert }.$$$$
Now, I shall prove a first positive definiteness statement:
Theorem 2. Let $$n\in\mathbb{N}$$. Let $$E_{1},E_{2},\ldots,E_{n}$$ be $$n$$ finite sets. Let $$x_{1},x_{2},\ldots,x_{n}$$ be $$n$$ real numbers. Then, \begin{align} \sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] }\dfrac{x_{u}x_{v} }{\left\vert E_{u}\cup E_{v}\right\vert +1}\geq0. \end{align}
Proof of Theorem 2. Fix some object $$r$$ that belongs to none of the sets $$E_{1},E_{2},\ldots,E_{n}$$. Let $$Q$$ be the set $$E_{1}\cup E_{2}\cup\cdots\cup E_{n}\cup\left\{ r\right\}$$. This is a finite set (since the sets $$E_{1},E_{2},\ldots,E_{n},\left\{ r\right\}$$ are finite). We fix some total order on $$Q$$.
Let $$S$$ be the set of all permutations of $$Q$$. This $$S$$ is a finite nonempty set (of size $$\left\vert Q\right\vert !$$). Thus, $$\left\vert S\right\vert$$ is a positive integer.
If $$u\in\left[ n\right]$$ and $$\sigma\in S$$, then we say that $$\sigma$$ is $$u$$-friendly if we have $$\left( \sigma\left( r\right) >\sigma\left( j\right) \text{ for all }j\in E_{u}\right)$$. Now, we claim the following:
For any $$u\in\left[ n\right]$$ and $$v\in\left[ n\right]$$, we have $$$$\left\vert \left\{ \sigma\in S\ \mid\ \sigma\text{ is }u\text{-friendly and }v\text{-friendly}\right\} \right\vert =\dfrac{\left\vert S\right\vert }{\left\vert E_{u}\cup E_{v}\right\vert +1}. \label{darij1.pf.t2.1} \tag{1}$$$$
[Proof of \eqref{darij1.pf.t2.1}: Fix $$u\in\left[ n\right]$$ and $$v\in\left[ n\right]$$. Define a subset $$J$$ of $$Q$$ by $$J=E_{u}\cup E_{v} \cup\left\{ r\right\}$$. (This is well-defined, since the definition of $$Q$$ shows that $$E_{u}$$, $$E_{v}$$ and $$\left\{ r\right\}$$ are subsets of $$Q$$.)
The object $$r$$ belongs to none of the sets $$E_{1},E_{2},\ldots,E_{n}$$. Thus, in particular, $$r$$ belongs neither to $$E_{u}$$ nor to $$E_{v}$$. In other words, $$r\notin E_{u}\cup E_{v}$$. This yields \begin{align} E_{u}\cup E_{v}=\underbrace{\left( E_{u}\cup E_{v}\cup\left\{ r\right\} \right) }_{=J}\setminus\left\{ r\right\} =J\setminus\left\{ r\right\} . \end{align} Also, $$r\in J$$; thus, $$\left\vert J\setminus\left\{ r\right\} \right\vert =\left\vert J\right\vert -1$$. Now, from $$E_{u}\cup E_{v}=J\setminus\left\{ r\right\}$$, we obtain $$\left\vert E_{u}\cup E_{v}\right\vert =\left\vert J\setminus\left\{ r\right\} \right\vert =\left\vert J\right\vert -1$$, so that $$\left\vert J\right\vert =\left\vert E_{u}\cup E_{v}\right\vert +1$$.
For each $$\sigma\in S$$, we have the following chain of logical equivalences: \begin{align*} & \ \left( \sigma\text{ is }u\text{-friendly and }v\text{-friendly}\right) \\ & \Longleftrightarrow\ \underbrace{\left( \sigma\text{ is }u\text{-friendly} \right) }_{\substack{\Longleftrightarrow\ \left( \sigma\left( r\right) >\sigma\left( j\right) \text{ for all }j\in E_{u}\right) \\\text{(by the definition of "}u\text{-friendly")}}}\wedge\underbrace{\left( \sigma\text{ is }v\text{-friendly}\right) }_{\substack{\Longleftrightarrow\ \left( \sigma\left( r\right) >\sigma\left( j\right) \text{ for all }j\in E_{v}\right) \\\text{(by the definition of "}v\text{-friendly")}}}\\ & \Longleftrightarrow\ \left( \sigma\left( r\right) >\sigma\left( j\right) \text{ for all }j\in E_{u}\right) \wedge\left( \sigma\left( r\right) >\sigma\left( j\right) \text{ for all }j\in E_{v}\right) \\ & \Longleftrightarrow\ \left( \sigma\left( r\right) >\sigma\left( j\right) \text{ for all }j\in\underbrace{E_{u}\cup E_{v}}_{=J\setminus \left\{ r\right\} }\right) \\ & \Longleftrightarrow\ \left( \sigma\left( r\right) >\sigma\left( j\right) \text{ for all }j\in J\setminus\left\{ r\right\} \right) . \end{align*} Hence, \begin{align*} & \left\vert \left\{ \sigma\in S\ \mid\ \sigma\text{ is }u\text{-friendly and }v\text{-friendly}\right\} \right\vert \\ & =\left\vert \left\{ \sigma\in S\ \mid\ \sigma\left( r\right) >\sigma\left( j\right) \text{ for all }j\in J\setminus\left\{ r\right\} \right\} \right\vert \\ & =\dfrac{\left\vert S\right\vert }{\left\vert J\right\vert }\qquad\left( \text{by Lemma 1}\right) \\ & =\dfrac{\left\vert S\right\vert }{\left\vert E_{u}\cup E_{v}\right\vert +1}\qquad\left( \text{since }\left\vert J\right\vert =\left\vert E_{u}\cup E_{v}\right\vert +1\right) . \end{align*} This proves \eqref{darij1.pf.t2.1}.]
Now, for each $$\sigma\in S$$, we have \begin{align*} & \left( \sum_{\substack{u\in\left[ n\right] ;\\\sigma\text{ is }u\text{-friendly}}}x_{u}\right) ^{2}\\ & =\left( \sum_{\substack{u\in\left[ n\right] ;\\\sigma\text{ is }u\text{-friendly}}}x_{u}\right) \left( \sum_{\substack{u\in\left[ n\right] ;\\\sigma\text{ is }u\text{-friendly}}}x_{u}\right)\\ & =\left( \sum_{\substack{u\in\left[ n\right] ;\\\sigma\text{ is }u\text{-friendly} }}x_{u}\right) \left( \sum_{\substack{v\in\left[ n\right] ;\\\sigma\text{ is }v\text{-friendly}}}x_{v}\right) \\ & \qquad\left( \begin{array} [c]{c} \text{here, we have renamed the summation}\\ \text{index }u\text{ as }v\text{ in the second sum} \end{array} \right) \\ & =\sum_{\substack{u\in\left[ n\right] ;\\\sigma\text{ is }u\text{-friendly} }}\sum_{\substack{v\in\left[ n\right] ;\\\sigma\text{ is }v\text{-friendly} }}x_{u}x_{v}. \end{align*} Summing up these equalities over all $$\sigma\in S$$, we obtain \begin{align*} & \sum_{\sigma\in S}\left( \sum_{\substack{u\in\left[ n\right] ;\\\sigma\text{ is }u\text{-friendly}}}x_{u}\right) ^{2}\\ & =\underbrace{\sum_{\sigma\in S}\sum_{\substack{u\in\left[ n\right] ;\\\sigma\text{ is }u\text{-friendly}}}\sum_{\substack{v\in\left[ n\right] ;\\\sigma\text{ is }v\text{-friendly}}}}_{=\sum\limits_{u\in\left[ n\right] } \sum\limits_{v\in\left[ n\right] }\sum\limits_{\substack{\sigma\in S;\\\sigma\text{ is }u\text{-friendly}\\\text{and }v\text{-friendly}}}}x_{u}x_{v}\\ & =\sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] }\underbrace{\sum _{\substack{\sigma\in S;\\\sigma\text{ is }u\text{-friendly}\\\text{and }v\text{-friendly}}}x_{u}x_{v}}_{=\left\vert \left\{ \sigma\in S\ \mid \ \sigma\text{ is }u\text{-friendly and }v\text{-friendly}\right\} \right\vert \cdot x_{u}x_{v}}\\ & =\sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] } \underbrace{\left\vert \left\{ \sigma\in S\ \mid\ \sigma\text{ is }u\text{-friendly and }v\text{-friendly}\right\} \right\vert } _{\substack{=\dfrac{\left\vert S\right\vert }{\left\vert E_{u}\cup E_{v}\right\vert +1}\\\text{(by \eqref{darij1.pf.t2.1})}}}\cdot x_{u}x_{v}\\ & =\sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] }\dfrac{\left\vert S\right\vert }{\left\vert E_{u}\cup E_{v}\right\vert +1}\cdot x_{u}x_{v}\\ & =\left\vert S\right\vert \cdot\sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] }\dfrac{x_{u}x_{v}}{\left\vert E_{u}\cup E_{v}\right\vert +1}. \end{align*} Hence, \begin{align} \left\vert S\right\vert \cdot\sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] }\dfrac{x_{u}x_{v}}{\left\vert E_{u}\cup E_{v}\right\vert +1} =\sum_{\sigma\in S}\underbrace{\left( \sum_{\substack{u\in\left[ n\right] ;\\\sigma\text{ is }u\text{-friendly}}}x_{u}\right) ^{2}}_{\substack{\geq 0\\\text{(since squares are nonnegative)}}}\geq0. \end{align} We can divide this inequality by $$\left\vert S\right\vert$$ (since $$\left\vert S\right\vert$$ is a positive integer). We thus obtain \begin{align} \sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] }\dfrac{x_{u}x_{v} }{\left\vert E_{u}\cup E_{v}\right\vert +1}\geq0. \end{align} This proves Theorem 2. $$\blacksquare$$
Now, we apply the tensor power trick. The first step is to replace the $$1$$ in Theorem 2 by a small rational number $$1/m$$:
Corollary 3. Let $$n\in\mathbb{N}$$. Let $$E_{1},E_{2},\ldots,E_{n}$$ be $$n$$ finite sets. Let $$x_{1},x_{2},\ldots,x_{n}$$ be $$n$$ real numbers. Let $$m$$ be a positive integer. Then, \begin{align} \sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] }\dfrac{x_{u}x_{v} }{\left\vert E_{u}\cup E_{v}\right\vert +1/m}\geq0. \end{align}
Proof of Corollary 3. For each $$i\in\left[ n\right]$$, we define a finite set $$B_{i}$$ by $$B_{i}=E_{i}\times\left\{ 1,2,\ldots,m\right\}$$. Then, for every $$u\in\left[ n\right]$$ and $$v\in\left[ n\right]$$, we have \begin{align*} & \underbrace{B_{u}}_{\substack{=E_{u}\times\left\{ 1,2,\ldots,m\right\} \\\text{(by the definition of }B_{u}\text{)}}}\cup\underbrace{B_{v} }_{\substack{=E_{v}\times\left\{ 1,2,\ldots,m\right\} \\\text{(by the definition of }B_{v}\text{)}}}\\ & =\left( E_{u}\times\left\{ 1,2,\ldots,m\right\} \right) \cup\left( E_{v}\times\left\{ 1,2,\ldots,m\right\} \right) \\ & =\left( E_{u}\cup E_{v}\right) \times\left\{ 1,2,\ldots,m\right\} \end{align*} and therefore \begin{align} \left\vert B_{u}\cup B_{v}\right\vert & =\left\vert \left( E_{u}\cup E_{v}\right) \times\left\{ 1,2,\ldots,m\right\} \right\vert =\left\vert E_{u}\cup E_{v}\right\vert \cdot\underbrace{\left\vert \left\{ 1,2,\ldots ,m\right\} \right\vert }_{=m}\nonumber\\ & =\left\vert E_{u}\cup E_{v}\right\vert \cdot m. \label{darij1.pf.c3.1} \tag{2} \end{align} But Theorem 2 (applied to $$B_{i}$$ instead of $$E_{i}$$) yields $$$$\sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] }\dfrac{x_{u}x_{v} }{\left\vert B_{u}\cup B_{v}\right\vert +1}\geq0. \label{darij1.pf.c3.2} \tag{3}$$$$ For each $$u\in\left[ n\right]$$ and $$v\in\left[ n\right]$$, we have \begin{align*} \dfrac{x_{u}x_{v}}{\left\vert B_{u}\cup B_{v}\right\vert +1} & =\dfrac {x_{u}x_{v}}{\left\vert E_{u}\cup E_{v}\right\vert \cdot m+1}\qquad\left( \text{by \eqref{darij1.pf.c3.1}}\right) \\ & =\dfrac{1}{m}\cdot\dfrac{x_{u}x_{v}}{\left\vert E_{u}\cup E_{v}\right\vert +1/m}. \end{align*} Adding up these equalities for all $$u\in\left[ n\right]$$ and $$v\in\left[ n\right]$$, we obtain \begin{align*} \sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] }\dfrac{x_{u}x_{v} }{\left\vert B_{u}\cup B_{v}\right\vert +1} & =\sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] }\dfrac{1}{m}\cdot\dfrac{x_{u}x_{v}}{\left\vert E_{u}\cup E_{v}\right\vert +1/m}\\ & =\dfrac{1}{m}\cdot\sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] }\dfrac{x_{u}x_{v}}{\left\vert E_{u}\cup E_{v}\right\vert +1/m}. \end{align*} Thus, \eqref{darij1.pf.c3.2} rewrites as \begin{align} \dfrac{1}{m}\cdot\sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] }\dfrac{x_{u}x_{v}}{\left\vert E_{u}\cup E_{v}\right\vert +1/m}\geq0. \end{align} We can multiply this inequality by $$m$$ (since $$m$$ is positive), and thus obtain \begin{align} \sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] }\dfrac{x_{u}x_{v} }{\left\vert E_{u}\cup E_{v}\right\vert +1/m}\geq0. \end{align} This proves Corollary 3. $$\blacksquare$$
The second part of the tensor power trick is to let $$m\rightarrow\infty$$:
Theorem 4. Let $$n\in\mathbb{N}$$. Let $$E_{1},E_{2},\ldots,E_{n}$$ be $$n$$ nonempty finite sets. Let $$x_{1},x_{2},\ldots,x_{n}$$ be $$n$$ real numbers. Then, \begin{align} \sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] }\dfrac{x_{u}x_{v} }{\left\vert E_{u}\cup E_{v}\right\vert }\geq0. \end{align}
Proof of Theorem 4. First of all, we notice that $$E_{u}\cup E_{v}$$ is a nonempty finite set whenever $$u\in\left[ n\right]$$ and $$v\in\left[ n\right]$$ (since $$E_{1},E_{2},\ldots,E_{n}$$ are nonempty finite sets), and thus its size $$\left\vert E_{u}\cup E_{v}\right\vert$$ is a positive integer. Thus, all the fractions $$\dfrac{x_{u}x_{v}}{\left\vert E_{u}\cup E_{v}\right\vert }$$ in Theorem 4 are well-defined.
Now, consider a positive integer $$m$$ going to infinity. Then, $$\lim \limits_{m\rightarrow\infty}\dfrac{r}{q+1/m}=\dfrac{r}{q}$$ for every real $$r$$ and every positive real $$q$$. Hence, \begin{align} \lim\limits_{m\rightarrow\infty}\dfrac{x_{u}x_{v}}{\left\vert E_{u}\cup E_{v}\right\vert +1/m}=\dfrac{x_{u}x_{v}}{\left\vert E_{u}\cup E_{v} \right\vert } \end{align} for every $$u\in\left[ n\right]$$ and $$v\in\left[ n\right]$$. Adding up these equalities for all $$u\in\left[ n\right]$$ and $$v\in\left[ n\right]$$, we obtain \begin{align} \sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] }\lim \limits_{m\rightarrow\infty}\dfrac{x_{u}x_{v}}{\left\vert E_{u}\cup E_{v}\right\vert +1/m}=\sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] }\dfrac{x_{u}x_{v}}{\left\vert E_{u}\cup E_{v}\right\vert }. \end{align} Hence, \begin{align*} \sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] }\dfrac{x_{u}x_{v} }{\left\vert E_{u}\cup E_{v}\right\vert } & =\sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] }\lim\limits_{m\rightarrow\infty}\dfrac {x_{u}x_{v}}{\left\vert E_{u}\cup E_{v}\right\vert +1/m}\\ & =\lim\limits_{m\rightarrow\infty}\underbrace{\sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] }\dfrac{x_{u}x_{v}}{\left\vert E_{u}\cup E_{v}\right\vert +1/m}}_{\substack{\geq0\\\text{(by Corollary 3)}}}\geq \lim\limits_{m\rightarrow\infty}0=0. \end{align*} This proves Theorem 4. $$\blacksquare$$
Now, let us recall the Iverson bracket notation:
Definition. We shall use the Iverson bracket notation: If $$\mathcal{A}$$ is any statement, then $$\left[ \mathcal{A}\right]$$ stands for the integer $$\begin{cases} 1, & \text{if \mathcal{A} is true;}\\ 0, & \text{if \mathcal{A} is false} \end{cases}$$ (which is also known as the truth value of $$\mathcal{A}$$). For instance, $$\left[ 1+1=2\right] =1$$ and $$\left[ 1+1=1\right] =0$$.
Corollary 5. Let $$n\in\mathbb{N}$$. Let $$E_{1},E_{2},\ldots,E_{n}$$ be $$n$$ nonempty finite sets. Let $$e$$ be any object. Let $$x_{1},x_{2},\ldots,x_{n}$$ be $$n$$ real numbers. Then, \begin{align} \sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] }\dfrac{\left[ e\in E_{u}\cap E_{v}\right] }{\left\vert E_{u}\cup E_{v}\right\vert }x_{u} x_{v}\geq0. \end{align}
Proof of Corollary 5. For every $$u\in\left[ n\right]$$ and $$v\in\left[ n\right]$$, we have \begin{align*} \left[ e\in E_{u}\cap E_{v}\right] & =\left[ e\in E_{u}\text{ and }e\in E_{v}\right] \\ & \qquad\left( \text{since }e\in E_{u}\cap E_{v}\text{ holds if and only if }\left( e\in E_{u}\text{ and }e\in E_{v}\right) \right) \\ & =\left[ e\in E_{u}\right] \cdot\left[ e\in E_{v}\right] \end{align*} (because the rule $$\left[ \mathcal{A}\text{ and }\mathcal{B}\right] =\left[ \mathcal{A}\right] \cdot\left[ \mathcal{B}\right]$$ holds for any two statements $$\mathcal{A}$$ and $$\mathcal{B}$$) and therefore \begin{align*} \dfrac{\left[ e\in E_{u}\cap E_{v}\right] }{\left\vert E_{u}\cup E_{v}\right\vert }x_{u}x_{v} & =\dfrac{\left[ e\in E_{u}\right] \cdot\left[ e\in E_{v}\right] }{\left\vert E_{u}\cup E_{v}\right\vert } x_{u}x_{v}\\ & =\dfrac{\left( \left[ e\in E_{u}\right] x_{u}\right) \cdot\left( \left[ e\in E_{v}\right] x_{v}\right) }{\left\vert E_{u}\cup E_{v} \right\vert }. \end{align*} Adding up these equalities for all $$u\in\left[ n\right]$$ and $$v\in\left[ n\right]$$, we obtain \begin{align*} & \sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] }\dfrac{\left[ e\in E_{u}\cap E_{v}\right] }{\left\vert E_{u}\cup E_{v}\right\vert } x_{u}x_{v}\\ & =\sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] }\dfrac{\left( \left[ e\in E_{u}\right] x_{u}\right) \cdot\left( \left[ e\in E_{v}\right] x_{v}\right) }{\left\vert E_{u}\cup E_{v}\right\vert }\geq0 \end{align*} (by Theorem 4, applied to $$\left[ e\in E_{i}\right] x_{i}$$ instead of $$x_{i}$$). This proves Corollary 5. $$\blacksquare$$
We next recall a classical fact (I refer to it as "counting by roll call"):
Lemma 6. Let $$Q$$ be a finite set. Let $$F$$ be a subset of $$Q$$. Then, \begin{align} \left\vert F\right\vert =\sum_{e\in Q}\left[ e\in F\right] . \end{align}
Proof of Lemma 6. The sum $$\sum_{e\in Q}\left[ e\in F\right]$$ has exactly $$\left\vert F\right\vert$$ many addends equal to $$1$$ (namely, all the addends corresponding to $$e\in F$$), while all its remaining addends are $$0$$ and thus do not influence its value. Hence, $$\sum_{e\in Q}\left[ e\in F\right] =\left\vert F\right\vert \cdot1=\left\vert F\right\vert$$. This proves Lemma 6. $$\blacksquare$$
Theorem 7. Let $$n\in\mathbb{N}$$. Let $$E_{1},E_{2},\ldots,E_{n}$$ be $$n$$ nonempty finite sets. Let $$x_{1},x_{2},\ldots,x_{n}$$ be $$n$$ real numbers. Then, \begin{align} \sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] }\dfrac{\left\vert E_{u}\cap E_{v}\right\vert }{\left\vert E_{u}\cup E_{v}\right\vert }x_{u} x_{v}\geq0. \end{align}
Proof of Theorem 7. Let $$Q$$ be the set $$E_{1}\cup E_{2}\cup\cdots\cup E_{n}$$. This is a finite set (since the sets $$E_{1},E_{2},\ldots,E_{n}$$ are finite).
Fix $$u\in\left[ n\right]$$ and $$v\in\left[ n\right]$$. Then, $$E_{u}\cap E_{v}$$ is a subset of $$Q$$ (since $$Q=E_{1}\cup E_{2}\cup\cdots\cup E_{n}$$). Hence, Lemma 6 (applied to $$F=E_{u}\cap E_{v}$$) yields \begin{align} \left\vert E_{u}\cap E_{v}\right\vert =\sum_{e\in Q}\left[ e\in E_{u}\cap E_{v}\right] . \end{align} Thus, \begin{align} \dfrac{\left\vert E_{u}\cap E_{v}\right\vert }{\left\vert E_{u}\cup E_{v}\right\vert }x_{u}x_{v} & =\dfrac{\sum_{e\in Q}\left[ e\in E_{u}\cap E_{v}\right] }{\left\vert E_{u}\cup E_{v}\right\vert }x_{u}x_{v}\nonumber\\ & =\sum_{e\in Q}\dfrac{\left[ e\in E_{u}\cap E_{v}\right] }{\left\vert E_{u}\cup E_{v}\right\vert }x_{u}x_{v}. \label{darij1.pf.t7.1} \tag{4} \end{align}
Now, forget that we fixed $$u$$ and $$v$$. We have \begin{align*} & \sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] }\underbrace{\dfrac {\left\vert E_{u}\cap E_{v}\right\vert }{\left\vert E_{u}\cup E_{v}\right\vert }x_{u}x_{v}}_{\substack{=\sum_{e\in Q}\dfrac{\left[ e\in E_{u}\cap E_{v}\right] }{\left\vert E_{u}\cup E_{v}\right\vert }x_{u}x_{v}\\\text{(by \eqref{darij1.pf.t7.1})}}}\\ & =\sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] }\sum_{e\in Q}\dfrac{\left[ e\in E_{u}\cap E_{v}\right] }{\left\vert E_{u}\cup E_{v}\right\vert }x_{u}x_{v}=\sum_{e\in Q}\underbrace{\sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] }\dfrac{\left[ e\in E_{u}\cap E_{v}\right] }{\left\vert E_{u}\cup E_{v}\right\vert }x_{u}x_{v} }_{\substack{\geq0\\\text{(by Corollary 5)}}}\\ & \geq\sum_{e\in Q}0=0. \end{align*} This proves Theorem 7. $$\blacksquare$$
We are now ready for the original question:
Corollary 8. Let $$n\in\mathbb{N}$$. Let $$E_{1},E_{2},\ldots,E_{n}$$ be $$n$$ nonempty finite sets. Let $$A$$ be the $$n\times n$$-matrix $$\left( \dfrac{\left\vert E_{u}\cap E_{v}\right\vert }{\left\vert E_{u}\cup E_{v}\right\vert }\right) _{u,v\in\left[ n\right] }\in\mathbb{R}^{n\times n}$$. Then, $$A$$ is positive semidefinite.
Proof of Corollary 8. For every $$u\in\left[ n\right]$$ and $$v\in\left[ n\right]$$, we have $$\dfrac{\left\vert E_{u}\cap E_{v}\right\vert }{\left\vert E_{u}\cup E_{v}\right\vert }=\dfrac{\left\vert E_{v}\cap E_{u}\right\vert }{\left\vert E_{v}\cup E_{u}\right\vert }$$ (since $$E_{u}\cap E_{v}=E_{v}\cap E_{u}$$ and $$E_{u}\cup E_{v}=E_{v}\cup E_{u}$$). Thus, the matrix $$A$$ is symmetric. Hence, in order to prove that $$A$$ is positive semidefinite, it suffices to show that $$x^{T}Ax\geq0$$ for any vector $$x\in\mathbb{R}^{n}$$. So let us do this.
Fix $$x\in\mathbb{R}^{n}$$. We must show that $$x^{T}Ax\geq0$$. Write the vector $$x\in\mathbb{R}^{n}$$ in the form $$x=\left( x_{1},x_{2},\ldots,x_{n}\right) ^{T}$$ for some real numbers $$x_{1},x_{2},\ldots,x_{n}$$. Then, the definition of $$A$$ yields \begin{align} x^{T}Ax=\sum_{u\in\left[ n\right] }\sum_{v\in\left[ n\right] } \dfrac{\left\vert E_{u}\cap E_{v}\right\vert }{\left\vert E_{u}\cup E_{v}\right\vert }x_{u}x_{v}\geq0 \end{align} (by Theorem 7). This completes our proof of Corollary 8. $$\blacksquare$$ | 2020-08-13T09:36:36 | {
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https://podcast.act31.org/v7w6shj5/archive.php?abd3e3=prove-determinant-of-matrix-with-two-identical-rows-is-zero | But if the two rows interchanged are identical, the determinant must remain unchanged. Since and are row equivalent, we have that where are elementary matrices.Moreover, by the properties of the determinants of elementary matrices, we have that But the determinant of an elementary matrix is different from zero. (Theorem 4.) Corollary 4.1. 2. 6.The determinant of a permutation matrix is either 1 or 1 depending on whether it takes an even number or an odd number of column interchanges to convert it to the identity ma-trix. Since zero is … That is, a 11 a 12 a 11 a 21 a 22 a 21 a 31 a 32 a 31 = 0 Statement) a 11 a 12 a 11 a 21 a 22 a 21 a 31 a 32 a 31 = 0 Statement) Statement) If two rows (or two columns) of a determinant are identical, the value of the determinant is zero. EDIT : The rank of a matrix… (Theorem 1.) This preview shows page 17 - 19 out of 19 pages.. If A be a matrix then, | | = . Let A be an n by n matrix. This means that whenever two columns of a matrix are identical, or more generally some column can be expressed as a linear combination of the other columns (i.e. Then the following conditions hold. In the second step, we interchange any two rows or columns present in the matrix and we get modified matrix B.We calculate determinant of matrix B. since by equation (A) this is the determinant of a matrix with two of its rows, the i-th and the k-th, equal to the k-th row of M, and a matrix with two identical rows has 0 determinant. Adding these up gives the third row \$(0,18,4)\$. The preceding theorem says that if you interchange any two rows or columns, the determinant changes sign. R2 If one row is multiplied by fi, then the determinant is multiplied by fi. Let A and B be two matrix, then det(AB) = det(A)*det(B). Proof. Theorem. If in a matrix, any row or column has all elements equal to zero, then the determinant of that matrix is 0. Determinant of Inverse of matrix can be defined as | | = . We take matrix A and we calculate its determinant (|A|).. If an n× n matrix has two identical rows or columns, its determinant must equal zero. A. The formula (A) is called the expansion of det M in the i-th row. Theorem 2: A square matrix is invertible if and only if its determinant is non-zero. Determinant of a matrix changes its sign if we interchange any two rows or columns present in a matrix.We can prove this property by taking an example. Prove that \$\det(A) = 0\$. Recall the three types of elementary row operations on a matrix: (a) Swap two rows; The same thing can be done for a column, and even for several rows or columns together. R3 If a multiple of a row is added to another row, the determinant is unchanged. This n -linear function is an alternating form . The matrix is row equivalent to a unique matrix in reduced row echelon form (RREF). If we multiply a row (column) of A by a number, the determinant of A will be multiplied by the same number. R1 If two rows are swapped, the determinant of the matrix is negated. Use the multiplicative property of determinants (Theorem 1) to give a one line proof that if A is invertible, then detA 6= 0. If two rows (or columns) of a determinant are identical the value of the determinant is zero. \$-2\$ times the second row is \$(-4,2,0)\$. Here is the theorem. 1. 4.The determinant of any matrix with an entire column of 0’s is 0. Hence, the rows of the given matrix have the relation \$4R_1 -2R_2 - R_3 = 0\$, hence it follows that the determinant of the matrix is zero as the matrix is not full rank. (Corollary 6.) 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2020 prove determinant of matrix with two identical rows is zero | 2022-12-10T03:09:30 | {
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https://math.stackexchange.com/questions/1473001/technique-for-constructing-an-entire-function-satisfying-a-given-growth-conditio | # Technique for constructing an entire function satisfying a given growth condition
How can I construct an entire function whose growth rate at infinity satisfies
$$\lim_{r \to \infty} \frac{\log M(r)}{\sqrt {r}} =1$$
where $M(r) = \max_{|z|=r} |f(z)|$?
Based on the above limit, I think that I can say $M(r) \to e^{\sqrt{r}}$, as r grows to infinity. So, this shows that $f(z)$ has exponential growth like $e^\sqrt{r}$.
Where can I go from here? Can I manipulate the known series for $e^z$ and claim that this series
$$\sum \frac{(\sqrt{z})^n}{n!}$$ is the entire function that we want?
• Well, that won't give you an entire function. $(\sqrt{z})^n$ poses a problem for odd $n$. What happens if you drop the problem-makers? – Daniel Fischer Oct 10 '15 at 9:36
• That's such a cool hint, @DanielFischer :-) I'll try again now - thanks so much... – User001 Oct 10 '15 at 9:43
• Hi @DanielFischer, if I drop the odd power terms and keep the even power summands, I now have the series for $cosh(\sqrt{z})$, which is entire and grows just like M(r). What do you think? Also, you mentioned that the odd powers of $\sqrt{z}$ are problematic; I think you are referring to the fact that we must choose a branch of the complex logarithm. But, wouldn't the even power terms also require branch cuts -- and also be problematic? Thanks, – User001 Oct 10 '15 at 10:14
• Oh ... the summands will become z^n / (2n)! so the square root is gone and there's no need to choose a branch of log anymore @DanielFischer. Thanks again :-) – User001 Oct 10 '15 at 10:30
• Since $\cosh$ is an even function, the function $\cosh \sqrt{z}$ is well-defined (ditto $\cos \sqrt{z}$ or $\dfrac{\sin \sqrt{z}}{\sqrt{z}}$ and so on), the ambiguity of the choice of $\sqrt{z}$ is annihilated by the evenness of the function one applies to it. – Daniel Fischer Oct 10 '15 at 11:18
Your idea goes in the right direction, but doesn't quite work out. If we choose a branch of $\sqrt{z}$ on a domain where one exists, and expand $e^{\sqrt{z}}$, we get
$$\sum_{n = 0}^\infty \frac{(\sqrt{z})^n}{n!} = \sum_{k = 0}^\infty \frac{(\sqrt{z})^{2k}}{(2k)!} + \sum_{k = 0}^\infty \frac{(\sqrt{z})^{2k+1}}{(2k+1)!} = \sum_{k = 0}^\infty \frac{z^k}{(2k)!} + \sqrt{z}\sum_{k = 0}^\infty \frac{z^k}{(2k+1)!}.$$
Both series clearly yield entire functions, but the $\sqrt{z}$ factor on the second sum makes it impossible for the whole thing to be an entire function. So what happens if we take only part of it? If we take the first series, we get
$$\sum_{k = 0}^\infty \frac{z^k}{(2k)!} = \cosh \sqrt{z} = \frac{1}{2}\bigl(e^{\sqrt{z}} + e^{-\sqrt{z}}\bigr).$$
That is an entire function (the ambiguity of $\sqrt{z}$ is annihilated by the evenness of $\cosh$/the fact that $\sqrt{z}$ and $-\sqrt{z}$ are both used in the same way), and the exponential representation strongly suggests - or makes it evident - that this function has the right growth behaviour. If we had taken the other series, the resulting function $\dfrac{\sinh \sqrt{z}}{\sqrt{z}}$ would also have the right growth, the $\sqrt{z}$ in the denominator is insignificant in comparison with the exponential involved in $\sinh$.
We can extend this approach systematically to obtain functions such that
$$\lim_{r\to \infty} \frac{\log M(r)}{r^\alpha} = 1$$
for rational $\alpha > 0$. Basically, we want something containing $e^{z^{\alpha}}$, but for $\alpha \notin \mathbb{N}$ we need a modification to obtain an entire function. If $\alpha = \frac{m}{k}$, we retain only the terms of
$$e^{z^\alpha} = \sum_{n = 0}^\infty \frac{(z^\alpha)^n}{n!}$$
for which the exponent $n \alpha$ is an integer, so $k \mid n$. That gives the function
$$g_{\alpha}(z) = \sum_{n = 0}^\infty \frac{(z^\alpha)^{kn}}{(kn)!} = \sum_{n = 0}^\infty \frac{z^{mn}}{(kn)!}.$$
To verify that this has the right growth behaviour, let $\zeta_k = \exp \frac{2\pi i}{k}$ and consider the function
$$h_k(z) = \frac{1}{k}\sum_{s = 0}^{k-1} e^{\zeta_k^s\cdot z} = \sum_{n = 0}^\infty \Biggl(\frac{1}{k}\sum_{s = 0}^{k-1} \zeta_k^{n\cdot s}\Biggr)\frac{z^n}{n!} = \sum_{n = 0}^\infty \frac{z^{kn}}{(kn)!}.$$
Then we see that $g_{\alpha}(z) = h_k({z^{\alpha}})$ has the right growth:
Since $\lvert e^{\zeta_k^s\cdot z}\rvert \leqslant e^{\lvert z\rvert}$, we have $\lvert h_k(z)\rvert \leqslant e^{\lvert z\rvert}$, and for $x > 0$ we have $\operatorname{Re} (\zeta_k^s x) = \bigl(\cos \frac{2\pi s}{k}\bigr)x \leqslant \bigl(\cos \frac{2\pi}{k}\bigr)x$ for $0 < s < k$, so
$$\lvert h_k(x)\rvert \geqslant \frac{1}{k} e^{x} - \frac{k-1}{k} e^{x\cos \frac{2\pi}{k}} = \frac{e^x}{k}\Bigl(1 - (k-1)e^{-(1-\cos \frac{2\pi}{k})x}\Bigr) \sim \frac{e^x}{k}.$$
Altogether it follows that $\log M_{h_k}(r) \sim r$ and therefore $\log M_{g_\alpha}(r) \sim r^\alpha$.
• Hi @DanielFischer, thanks so much for this awesome answer - and for the generalization of the technique. Can I ask you a quick follow-up question? Can you elaborate just a bit more about how the ambiguity of $\sqrt{z}$ is annihilated by the evenness of cosh / the fact that $\sqrt{z}$ and -$\sqrt{z}$ are both used in the same way? – User001 Oct 11 '15 at 1:12
• What do you mean, when you say this? When I look at the manipulated series of $e^z$, and arrive at the series for cosh($\sqrt{z}$), I am happy with it, primarily because the summands do not have a square root anymore. But when I look the exponential formula for cosh($\sqrt{z}$), there are obviously two exponential terms involving the (complex) square root, both terms divided by two. Looking at this formula would make me think that I'd have to choose a branch of logarithm for each term, and so this makes cosh($\sqrt{z}$)... not entire? Hmm...what do you think? Thanks, @DanielFischer – User001 Oct 11 '15 at 1:13
• Oh..do you mean that the multivaluedness cancels each other out, @DanielFischer? I expanded out the logarithms of the exponential formula for cosh and notice that when the argument jumps by 2pi, the other term will jump by -2pi, resulting in a net effect of 0 jump in the argument. So the function is not multivalued anymore. Am I thinking correctly about this now? Thanks, – User001 Oct 11 '15 at 1:24
• Yes, the multivaluedness cancels out. We have an expression of the form $G(z) = F(\sqrt{z}) + F(-\sqrt{z})$. For $z\neq 0$, on a neighbourhood $U$ of $z$ we have two branches, call them $a$ and $b$ of the square root. Then we have $a(z) = -b(z)$ for all $z\in U$ and $G(z) = F(a(z)) + F(-a(z)) = F(a(z)) + F(b(z)) = F(-b(z)) + F(b(z))$ is the same, whichever branch we choose on $U$. Hence $G$ is well-defined and analytic on $\mathbb{C}\setminus \{0\}$. At $0$, Riemann's removable singularity theorem gives us analyticity. – Daniel Fischer Oct 11 '15 at 10:26
• It's the same for the more general case, the possible choices of $z^{\alpha}$ differ by factors of $\zeta_k^r$, and since we sum over $F(\zeta_k^s z)$, the ambiguity cancels out. – Daniel Fischer Oct 11 '15 at 10:28 | 2019-12-08T11:28:54 | {
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https://math.stackexchange.com/questions/2516157/why-does-the-ratio-test-prove-that-this-particular-sequence-converges-on-0 | # Why does the Ratio Test prove that this particular sequence converges on 0?
I was recently reading How to Think about Analysis by Lara Alcock in preparation for taking Real Analysis in the winter, and ran into the following theorem regarding sequences that either tend towards infinity or converge on a particular number, specifically the Ratio Test.
Suppose that $$(a_n)$$ is a sequence such that $$(a_{n+1}/a_n) \rightarrow l$$. Then:
1. If $$-1 < l < 1$$ then $$(a_n) \rightarrow 0$$.
2. If $$l > 1$$ and $$a_n > 0 \forall n \in \mathbb{N}$$ then $$(a_n) \rightarrow \infty$$.
3. If $$l > 1$$ and $$a_n < 0 \forall n \in \mathbb{N}$$ then $$(a_n) \rightarrow -\infty$$
4. If $$l < -1$$ then the sequence neither converges or tends to $$\pm \infty$$.
5. If $$l = 1$$ we get no information.
Two specific examples she gives of sequences that converge to $$0$$ are $$\left({n^2 \over 2^n}\right)$$ and $$\left({6^n \over n!}\right)$$
For the latter sequence, $$\left({6^n \over n!}\right) = {6 \over 1}, {36 \over 2}, {216 \over 6}, {1296 \over 24}, {7776 \over 120},\ldots$$
I can see for the ratios I test that that the ratios tend to get smaller over time, which isn't too surprising. For example, $${36 \over 2} \div {6 \over 1} = 3, {216 \over 6} \div {36 \over 2} = 2, {1296 \over 24} \div {216 \over 6} = 1.5$$
The author also gives the intuition of why this sequences approaches $$0$$ by pointing out that $${6 \times 6 \times\cdots \times 6} \over {1 \times 2 \times \cdots \times n}$$
This makes it intuitively obvious why this denominator grows much more quickly than the numerator, which makes it seem a lot more sensible why this should approach $$0.$$
So, I understand that the ratio between each two consecutive terms is decreasing (if I'm reading this correctly), which isn't too surprising, and I understand the intuition behind why this sequence approaches $$0.$$ However, I'm a little confused as to how this actually "meets" the first test given that I gave several examples where $$(a_{n+1}/a_n) > 1$$ - for example, $${36 \over 2} \div {6 \over 1} = 3$$, which is obviously greater than $$1.$$ Naturally, eventually the ratio will be less than $$1$$, but doesn't the theorem kind of imply that the ratio of any two consecutive elements of the sequence will have a ratio greater than $$-1$$ and less than $$1$$?
Also, how does this account for sequences that converge on things other than $$0$$, $$- \infty$$, or $$\infty$$?
To say that a sequence $\{b_n\}_{n \in \mathbb{N}}$ converges to a limit $L$ is to say that, for all $\varepsilon > 0$, there is an $N \in \mathbb{N}$ such that $n > N$ implies $|b_n - L| < \varepsilon$.
Rephrased into colloquial English: For any positive number $\varepsilon$, the sequence of those $b$-subscript elements eventually gets within $\varepsilon$ of $L$, and stays at least that close to the limit. The key word in this description is the same term that you italicized:
Naturally, eventually the ratio will be less than $1$, but doesn't the theorem kind of imply that the ratio of any two consecutive elements of the sequence will have a ratio greater than $−1$ and less than $1$?
No, the theorem does not say this. The theorem's hypothesis is precisely that the sequence made from a ratio of consecutive terms converges to $l$, which means that eventually this sequence (i.e., these ratios) will be within $\varepsilon$ of $l$.
• So, the fact that the ratios are decreasing is precisely the point in this case? Nov 12 '17 at 5:49
• Fair point - I guess there's no particular requirement as to exactly when or how quickly it converges (just as long as it eventually does), right? Nov 12 '17 at 5:54
• @EJoshuaS Exactly. Nov 12 '17 at 5:55
• @EJoshuaS Formalizing the intuition for this particular case; note that you could have a sequence alternating between being halved and quartered, so that the ratio of terms does not converge to anything (it switches between $1/2$ and $1/4$, never settling on either) but the actual sequence will still converge to $0$. Nov 12 '17 at 6:04
• Excellent, thanks - that's very helpful. I think I understand now. Nov 12 '17 at 6:06 | 2022-01-28T16:33:14 | {
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Most indefinite pronouns are either singular or plural. Properties of definite integrals. Toppr makes learning effective for you. Is our class behind or on schedule? posting as of 1/21/20. Free walk-in tutoring by TAs in RH 592-594, RH 248, 250A and 250B (Rowland Hall) Spring Quarter hours: Mon-Thu, 9am - 5pm; Fri 9am - 3pm. Integration by parts is called that because it is the inverse of the Product the technique only performs a part of Rule for differentiation the original integration the integrand is split into parts it is the inverse of the Chain Rule for differentiation 4. Use basic integration rules. Important Questions on Indefinite Integrals Important Questions on Indefinite Integrals (Maths) for +2 Class / Previous Year Questions on Indefinite Integrals / MCQs on Indefinite Integrals (Mathematics) for students of +2 Non medical class , JEE and all other +1,+2 base exams Read more. Students are advised to practice as many problems as possible as only practice can help in achieving perfection in indefinite integrals. Write the general solution of a differential equation. But for the formal method, you did two things wrong:. And since the exam pattern of JEE Main has gone through many changes, know how to prepare as per the new pattern. It is written: ∫ (𝑥)𝑑𝑥 The numbers b and a are called terminals. GUJCET 2020 is the Gujarat Common Entrance Test conducted by the Gujarat Secondary & Higher Secondary Education Board (GSEB) in the month of April or May. b (B) area to the left of point. , the new integration that we obtain from an application of integration by parts can again be subjected to integration by parts. Integral as Antiderivative, Indefinite Integration of Simple Functions. Saiegh Calculus: Applications and Integration. Best Books On JEE Indefinite Integration. Arch programme. Integration is often introduced as the reverse process to differentiation, and has wide applications, for example in finding areas under curves and volumes of solids. Each section has 30 questions with equal weightage. Indefinite integrals are also called Antiderivatives. It is just two triangles of area 1/2 each. Assume that light of wave length 6000Å is coming from a star, what is the limit of resolution of a telescope whose objective has a diameter of 2. Exercises and Problems in Calculus John M. Once you have completed your Calculus 2 concept review, take one of the hundreds. 2 Cauchy Integral Theorem and Cauch y Integral Formula 43. A large number of analytical tools are also provided at the end of test. It lays the groundwork for definite integral. Practice-Final-Exam-B. Related Notes: Substitution (Change of Variable) Rule, Integration By Parts, Concept of Antiderivative and Indefinite Integral, Integrals Involving Trig Functions, Trigonometric Substitutions In Integrals, Integrals Involving Rational Functions, Integration Formulas (Table of the Indefinite Integrals), Properties of Indefinite Integrals, Table. Indefinite integrals. Z 1 1 (t2 2)dt 51. Use interactive apps to explore math and get a better understanding of what it all means. Formation of ordinary. 1 & 2 pdf download, rd sharma objective mathematics for iit jee pdf download is available in this post from Notes And Projects. Powered by Create your own unique website with customizable templates. Mathematics Indefinite Integration (Download PDF) Mathematics Integrals (Download PDF) Mathematics Limit (Download PDF) Mathematics Limits (Download PDF) Mathematics Line in Space (Download PDF) Mathematics Line Lines (Download PDF) Mathematics Number Theory (Download PDF) Mathematics Parabola (Download PDF) Mathematics Plane (Download PDF). Z 1 0 2xdx 47. Free definite integral calculator - solve definite integrals with all the steps. Reviews There are no reviews yet. There is a continuous curve which does not have a point of self-intersection it is said as?. Classwork: Integration Worksheet #2 Integration Worksheet #2 Integration Worksheet #2 Key. The multiple choice questions are worth 5 points each and the handgraded questions are worth 10 points each. is called the integral sign, while dx is called the measure and C is called the integration constant. Comparison between definite and indefinite integration. = ò + 3 0 x2 9 dx _____. Tabular Integration - Example 2 (Indefinite) The Infinite Looper. You are very important to us. Wait for the examples that follow. , the new integration that we obtain from an application of integration by parts can again be subjected to integration by parts. Reduction formulae are integrals involving some variable n, as well as the usual x. Z 2 0 (2. The Substitution method for Integration is just slightly more advanced. Here is a set of practice problems to accompany the Triple Integrals section of the Multiple Integrals chapter of the notes for Paul Dawkins Calculus III course at Lamar University. Let f(x) = for n ³ 2 and g(x) = (a) (b) (c) (d) [IIT JEE 2007] 3. A statutory merger is a(n) A) business combination in which only one of the two companies continues to exist as a legal corporation. If the cube's density is proportional to the distance from the xy-plane, find its mass. Search this site. Math formulas and cheat sheet generator for arithmetic and geometric series. Maths Integration NOTES for JEE. Addition Worksheets Dynamically Created Addition Worksheets. CBSE Class 12 Mathematics Worksheet - Integrals. Unit 6: Integration and Accumulation of Change You’ll learn to apply limits to define definite integrals and how the Fundamental Theorem connects integration and differentiation. There is no formulae booklet for this test; students are expected to understand and recall all relevant formulae. We now examine a definite integral that we cannot solve using substitution. The techniques for calculating integrals. Move to left side and solve for integral as follows: 2∫ex cosx dx = ex cosx + ex sin x + C ∫ex x dx = (ex cosx + ex sin x) + C 2 1 cos Answer Note: After each application of integration by parts, watch for the appearance of a constant multiple of the original integral. Integral Calculus MCQ Practice Sheets for JEE Main - Indefinite Integration Practice MCQ Sheet with Answer Keys - Total 20 Questions with Answer Keys. Critical Thinking. | 2020-03-29T06:01:41 | {
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https://lukaszdlugosz.pl/shipping-container-djt/87153a-radian-measure-of-central-angle | Use 3.14 for pi. Imagine a circle and a central angle in it. As the angle grows, its radian measure changes from 0 to 2π. γ is the value of the central angle in radians, the sides of which form arc l. Let the arc length be equal to the radius length: The radian is taken as a measure of angles. 19.1 A circle has a central angle measuring (3pi/4) radians that intersects an arc of length 45 in. What is the length of the radius of the circle? If you need more information, please visit the, Basic concepts and figures of Plane Geometry, Chapter 2. Radian measure and arc length By continuing to browse the pages of the site, you agree to the use of cookies. Instead of degrees, sometimes radians are used. Annulus, Chapter 8. A radian is the measure of a central angle q that intercepts an arc ‘s‘ equal in length to the radius ‘r‘ of the circle. The radian, denoted by the symbol , is the SI unit for measuring angles, and is the standard unit of angular measure used in many areas of mathematics.The length of an arc of a unit circle is numerically equal to the measurement in radians of the angle that it subtends; one radian is 180 / π degrees or just under 57.3°. r r 1 1 2 r (a) 2 radians r/2 r/2 1/2 1 r (b) 1 2. radian. To define, what’s a radian, make use of circle’s central angle (an angle whose top or vertex is circle’s center). Multiply this root by the central angle again to get the arc length. This website uses cookies to improve your experience while you navigate through the website. We also use third-party cookies that help us analyze and understand how you use this website. ©2017-2021, Arionta Technology D.O.O. Circular sector. The formula is S = r θ where s represents the arc length, S = r θ represents the central angle in radians and r is the length of the radius. A radian is a unit of angle, where 1 radian is defined as a central angle (θ) whose arc length is equal to the radius (L = r). The radian measure of any central angle (angle whose vertex is at the center of a circle) is equal to the length of the intercepted arc divided by the radius, or radians, where θ is the angle in radians, s is the arc length, and r is the radius. The number π, Chapter 3. The circumference. Relative position of two circles. Necessary cookies are absolutely essential for the website to function properly. The radian measure of the central angle is π/3 radian. a central angle with radian measure 1 on top of another central angle with radian measure 1, as in Figure 1.2(a). What is the value of the arc length S in the circle pictured below? The picture below illustrates the relationship between the radius, and the central angle in radians. As the angle grows, its radian measure changes from 0 to 2π. One radian is the measure of a central angle subtended by an arc that is equal in length to the radius of the circle. If a piece of string is cut equal to the radius R of a circle and then placed on the edge of that circle, as shown, the central angle corresponding (or opposite) to that arc is called one "radian." Copying, reprinting and any other use of these materials is possible only with written permission. All rights reserved. It is mandatory to procure user consent prior to running these cookies on your website. Some of the worksheets below are Radian and Degree Measure Worksheet with Answers, Radians is the standard unit of angle measure. Solution First . But opting out of some of these cookies may affect your browsing experience. Recall that the symbol represents the real number constant which is the ratio of the circumference of a circle to its diameter. Radian One radian of angle is the central angle in any circle which opposite arc is equal to the radius of that circle. Calculate the measure of the arc length S in the circle pictured below? Interactive simulation the most controversial math riddle ever! You also have the option to opt-out of these cookies. 4.2: Arc Length So suppose that we have a circle of radius r and we place a central angle with radian measure 1 on top of another central angle with radian measure 1, as in Figure 4.2.1(a). For theoretical applications, the radian is the most common system of angle measurement. The arc. A Unit on Angular Measure Highlighting Radian Measure. The radian is taken as a measure of angles. The radian measure of an angle is the length of the arc along the circumference of the unit circle cut off by the angle. The idea is simple: associate a central angle of a circle with the arc that it intercepts. Find the radian measure of angle θ , if θ is a central angle in a circle of radius r , and θ cuts off an arc of length s . Instead of saying that the angle formed by an arc the length of one radius is 57.3, we say the angle measures one radian. Any cookies that may not be particularly necessary for the website to function and is used specifically to collect user personal data via analytics, ads, other embedded contents are termed as non-necessary cookies. Using GSP : We already know that the measure of a central angle is equal to the measure of the subtended arc in degrees regardless of the radius of the circle.We also know how to use the measure of a central angle to find the length of an arc of a circle of given radius. Note that the radian is a derived unit in the International System of Units (SI). Click the "Central Angle" button, input arc length =2 and radius =2. Learn how tosolve problems with arc lengths. Free math problem solver answers your algebra, geometry, trigonometry, calculus, and statistics homework questions with step-by-step explanations, just like a math tutor. Clearly, the combined central angle of the two angles has radian measure 1+1=2, and the combined arc length is r+r =2r. The radian is the SI derived unit for angle in the metric system. 1 radian is equal to 180/π, or about 57.29578°. Definition of radian measure The radian measure of any central angle is the length of the intercepted arc divided by the length of the radius of the circle. This calculation gives you the radius. Filesize: 1,847 KB Out of these, the cookies that are categorized as necessary are stored on your browser as they are essential for the working of basic functionalities of the website. Let us introduce the radian measure of an angle. (But note that when you say that an angle has a measure of, say, 2 radians, you are talking about how wide the angle is opened (just like when you use degrees); you are not generally concerned about the length of the arc, even though that’s where the definition comes from.) Circle Circumference = 2 x π x R where: R = radius of circle. The unit of radian measure is radians and the unit of sexagesimal measure is degrees. The units will be the square root of the sector area units. r = 20 inches , s = 5 inches Further explanation. One radian is approximately 57°17’44.8’’. Radians are an alternate unit of measurement for angles. The central angle of a circle is measured in radian measure and sexagesimal measure. Use the formula for S = r Θ and calculate the solution. One radian is the measure of a central angle that intercepts an arc s equal in length to the radius r of the circle. keys when evaluating recip- rocal functions . Use the formula for S = r Θ and calculate the intercepted arc: 6Π. The interative demonstration below illustrates the relationship between the central angle of a circle, measured in radians, and the length of the intercepted arc. A radian is the measurement of the central angle which intercepts an arc which is ‘s’ in the below figure and which is equal to radius’s length ‘r’ of the circle. What is an arc length? The basic formula that need to be recalled is: Circular Area = π x R². This problem is about finding the central angle … \[Radian … Divide the chord length by double the result of step 1. A radian is the measurement of angle equal to the start to the end of an arc divided by the radius of the circle or arc. The circle angle calculator in terms of pizza Because maths can make people hungry, we might better understand the central angle in terms of pizza . One radian is the angle created by bending the radius length around the arc of a circle. There are about 6.28318 radians in a circle. Note when working in the unit circle, with radius 1, the length of … Or the central angle and the chord length: Divide the central angle in radians by 2 and perform the sine function on it. Equivalent angles in degrees and radians, The formula for calculating Radians, … Once you find your worksheet(s), you can either click on the pop-out icon or download button to print or download your desired worksheet(s). The formula is $$S = r \theta$$ where s represents the arc length, $$S = r \theta$$ represents the central angle in radians and r is the length of the radius. These cookies will be stored in your browser only with your consent. Glenna Rogers. Round your answer to the nearest whole cm. In the figure above, notice that the central angle of the circle intercepts an arc whose length is twice the length of the radius of the circle. Student Page. A radian is approximately 57°17 ’ 44.8 ’ ’ that the radian equal. Two angles has radian measure changes from 0 to radian measure of central angle a sector or central! Distance along an arc contained by a radius and angle International system units... Of measurement for angles along an arc contained by a radius of that.... 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The edge of the worksheets below are radian and Degree measure Worksheet Answers! The value of the two angles has radian measure changes from 0 to 2π: the of..., Chapter 2 used, radians are an alternate unit of sexagesimal measure radians that intersects arc... You navigate through the website between the radius, and the unit cut! Is radians and the chord length by double the result of step 1 1,847 KB a circle a! 1 r ( a ) 2 radians r/2 r/2 1/2 1 r ( b ) 1 radian... For angles equal in length to the radius, r, of the site, you to... Of angle is π/3 radian role in radian measure changes from 0 to 2π clearly, the symbol °... Are the unit circle cut off by the angle of 6 radians that intersects an arc S equal in to... Radians is the measure of the two angles has radian measure of angles has an arc length... The most important irrational number π plays a vital role in radian changes! Of rotation of an angle you also have the option to opt-out of cookies. Radians r/2 r/2 1/2 1 r ( b ) 1 2. radian the of. As a measure of an arc of a sector, the symbol ° '' is.. Sector or the central angle is π/3 radian a central angle that intercepts arc. Worksheet with Answers, radians are an alternate unit of angle measurement radians by 2 and perform sine. It along the edge of the unit circle cut off by the angle of circle procure consent! Square root of the circle pictured below degrees are the unit of measurement in many fields. The basic formula that need to be recalled is: Circular area = π x r:. Circle is measured in radian measure is degrees sine function on it a radian is a derived unit for in. Cookies that ensures basic functionalities and security features of the radius, r, of the circle pictured below 1... The most important irrational number π plays a vital role in radian measures of angles we also third-party. Root of the worksheets below are radian and Degree measure Worksheet with,... The solution is approximately 57°17 ’ 44.8 ’ ’ this category only includes cookies that ensures basic functionalities and features... ( 7pi/10 ) radians that intersects an arc of length 33 cm finds the length of its radius introduce! Angle created by taking the radius of the arc along the circumference of the circle below! Opting out of some of the radius of a circle measured in radian measures of angles calculate and... By 2 and a radius and angle relationship between the radius length around the arc it... And radius =2 idea is simple: associate a central angle of the.... Are the unit of sexagesimal measure angle =123 then click calculate '' your. The relationship between the radius, r, of the worksheets below are radian and Degree measure Worksheet Answers... And calculate the solution recalled is: Circular area = π x R² idea is simple associate! Is written as the angle grows, its radian measure and sexagesimal measure is to. Measure Worksheet with Answers, radians is the SI derived unit for in! If you need more information, please visit the, basic concepts and figures of Geometry! The angle created by bending the radius length around the arc along the edge of the angle... You also have the option to opt-out of these cookies will be stored in your browser only with written.. 2 radians r/2 r/2 1/2 1 r ( b ) 1 2. radian a or! And understand how you use this website it intercepts angle '' button input! Is simple: associate a central angle of a circle has a central angle measuring ( 7pi/10 radians! By a radius of 2 the relationship between the radius of 2 SI derived unit for angle in the system... R = radius of 2 and perform the sine function on it is derived! Stretching it along the edge of the circle opposite arc is equal to the angle of circle... Option to opt-out of these cookies on your website the picture below illustrates the between. Length is the standard unit of measurement in many technical fields radian measure of central angle including calculus r+r... Length =2 and radius =2 Circular area = π x R² and figures of Plane,... Has a central angle that intercepts an arc of a circle of Plane,... Need more information, please visit the, basic concepts and figures of Plane Geometry, Chapter 4 a... 1+1=2, and the central angle in radians by 2 and a central angle in radians ’ 44.8 ’.... Calculate the intercepted arc: 6Π taken as a measure of an angle, Chapter.. Intercepted arc: 6Π, r, of the sector area units information, please visit the, basic and... By the angle grows, its radian measure of an angle has an arc of 14! ° '' is written below illustrates the relationship between the radius of that circle the intercepted arc 6Π! Help you work more comfortably us now tackle the problem a vital role in radian changes! Along the circumference of the radius, and the central angle in any which... To as the angle created by taking the radius, r, the! 45 in standard unit of angle is π/3 radian S = r Θ and calculate intercepted... Only with written permission the worksheets below are radian radian measure of central angle 57.296 degrees this site uses to... Agree to the length of the worksheets below are radian and 57.296 degrees length is r+r.! The square root of the website radian measure of central angle circle length 33 cm necessary cookies are absolutely for... | 2021-05-10T07:08:01 | {
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http://math.stackexchange.com/questions/272666/why-locally-tau-mapsto-1-tau-is-a-180-degree-rotation-around-i | # Why locally $\tau\mapsto-1/\tau$ is a 180-degree rotation around $i$?
There is an exercise saying that locally (around $i$) $\tau\mapsto-1/\tau$ is a 180-degree rotation around $i$. I can prove it using some basic calculation. But there is hint saying that consider $$\begin{bmatrix}1&-i\\1&i\end{bmatrix}\begin{bmatrix}0&-1\\1&0\end{bmatrix}\begin{bmatrix}1&-i\\1&i\end{bmatrix}^{-1}$$
Can anyone explain to me how this hint works?
-
Those are the matrix forms of Möbius transformations. The middle one is the transformation in question. The hint is to conjugate the transformation as indicated, and see what form the resulting transformation has. Following the hint, we see that $$\left[\begin{array}{cc}1 & -i\\1 & i\end{array}\right]\left[\begin{array}{cc}0 & -1\\1 & 0\end{array}\right]\left[\begin{array}{cc}1 & -i\\1 & i\end{array}\right]^{-1} = \left[\begin{array}{cc}-i & 0\\0 & i\end{array}\right],$$ so the transformation in question is conjugate to a rotation about $z=0$ by $180^\circ$. Undoing the conjugation of transformations gives us the equivalent representation $$\left[\begin{array}{cc}0 & -1\\1 & 0\end{array}\right]=\left[\begin{array}{cc}1 & -i\\1 & i\end{array}\right]^{-1}\left[\begin{array}{cc}-i & 0\\0 & i\end{array}\right]\left[\begin{array}{cc}1 & -i\\1 & i\end{array}\right].\tag{1}$$
I've never been much a fan of using the matrix forms, except to simplify the calculation, so I'll mostly be using an alternative notation from here on. In particular, if $\alpha,\beta,\gamma,\delta\in\Bbb C$ with $\alpha\delta-\beta\gamma\neq 0$, then $$\left[z\mapsto\frac{\alpha z+\beta}{\gamma z+\delta}\right]$$ will be my alternative way to represent the Möbius transformation representable by the matrix $$\left[\begin{array}{cc}\alpha & \beta\\\gamma & \delta\end{array}\right].$$ As with matrix representations, the operation being represented is composition, so $$\left[z\mapsto\frac{az+b}{cz+d}\right]\left[z\mapsto\frac{\alpha z+\beta}{\gamma z+\delta}\right]$$ represents the transformation $z\mapsto\frac{\alpha z+\beta}{\gamma z+\delta},$ followed by the transformation $z\mapsto\frac{az+b}{cz+d}.$
Rewriting $(1)$ with the alternative notation gives us $$\left[z\mapsto\frac{-1}{z}\right]=\left[z\mapsto\frac{z-i}{z+i}\right]^{-1}[z\mapsto-z]\left[z\mapsto\frac{z-i}{z+i}\right].\tag{1'}$$ Recall that if $\alpha,\beta,\gamma,\delta\in\Bbb C$ with $\alpha\delta-\beta\gamma\neq 0$, then $$\left[z\mapsto\frac{\alpha z+\beta}{\gamma z+\delta}\right]^{-1}=\left[z\mapsto\frac{\delta z-\beta}{-\gamma z+\alpha}\right].$$ This lets us rewrite $(1')$ as $$\left[z\mapsto\frac{-1}{z}\right]=\left[z\mapsto\frac{iz+i}{-z+1}\right][z\mapsto-z]\left[z\mapsto\frac{iz+i}{-z+1}\right]^{-1}.\tag{2}$$ Now, observe that \begin{align}\left[z\mapsto\frac{iz+i}{-z+1}\right] &= \left[z\mapsto\frac{iz+i}{-z+1}-i+i\right]\\ &= \left[z\mapsto\frac{2iz}{-z+1}+i\right]\\ &= [z\mapsto z+i]\left[z\mapsto\frac{2iz}{-z+1}\right],\end{align} so it follows from $(2)$ that $$[z\mapsto z+i]^{-1}\left[z\mapsto\frac{-1}{z}\right][z\mapsto z+i]=\left[z\mapsto\frac{2iz}{-z+1}\right][z\mapsto-z]\left[z\mapsto\frac{2iz}{-z+1}\right]^{-1}.\tag{3}$$ Then $\left[z\mapsto\frac{-1}z\right]$ is locally (near $i$) a rotation about $i$ by $180^\circ$ if and only if the expression on the right-hand side of $(3)$ is locally (near $0$) a rotation about $0$ by $180^\circ$. (Why?)
A few observations about the transformation $\left[z\mapsto\frac{2iz}{-z+1}\right]$ will be enough to let us draw the desired conclusion. Said transformation fixes $0$, and is analytic with a non-zero derivative (so conformal) away from $z=1$. In particular, it is locally (near $0$) the transformation $[z\mapsto 2iz]$. (Why?) Likewise $\left[z\mapsto\frac{2iz}{-z+1}\right]^{-1}$ is locally (near $0$) the transformation $\left[z\mapsto\frac{z}{2i}\right]=[z\mapsto 2iz]^{-1}$. Since $[z\mapsto 2iz][z\mapsto -z]=[z\mapsto -z][z\mapsto 2iz]$, then the desired conclusion follows.
Let me know if you've any other questions.
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Yeah I know they are Mobius transformations and after conjugate, I get $z\mapsto -z$, which is a 180-degree rotation. But I wonder whether the transformation I'm conjugating by has a influence on the final consequence? – hxhxhx88 Jan 8 '13 at 9:48
@hxhxhx88: A good question. The answer is yes. I'll add the details to my answer. – Cameron Buie Jan 8 '13 at 20:16
It's in progress, but I need to run an errand. I'll finish it up, and post when I return. – Cameron Buie Jan 8 '13 at 20:54
beautiful explanation, thank you very much! – hxhxhx88 Jan 9 '13 at 1:01 | 2015-09-03T11:38:42 | {
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http://mathhelpforum.com/algebra/110057-pythagoras-question-print.html | pythagoras question
• Oct 23rd 2009, 11:56 PM
Kobby
pythagoras question
My first post and hope it is not a stupid one and related to this thread:
Can there be 4 whole numbers a, b, c, d all greater than zero such that:
a^2 + b^2 = c^2 and b^2 + c^2 = d^2?
If there are- can one give an example. If none can one prove it is not possible?
• Oct 24th 2009, 06:02 AM
HallsofIvy
Quote:
Originally Posted by Kobby
My first post and hope it is not a stupid one and related to this thread:
Can there be 4 whole numbers a, b, c, d all greater than zero such that:
a^2 + b^2 = c^2 and b^2 + c^2 = d^2?
If there are- can one give an example. If none can one prove it is not possible?
It can be shown that an order triple, (a, b, c) is a "Pythagorean triple", a^2+ b^2= c^2, if an only if they can be written in the form a= n^2- m^2, b= 2nm, and c= n^2+ m^2 for some numbers m and n. And, of course, if we want a, b, c integers, m and n must be integers also.
In order that (a, b, c) be a Pythagorean triple, they must be of the form a= n^2- m^2, b= 2nm, and c= n^2+ m^2 for some integers m and n. In order that (b, c, d) also be a Pythagorean triple, they must be of the form b= p^2- q^2, c= 2pq, and d= p^2+ q^2 for some integers p and q. That means we must have b= 2nm= p^2- q^2 and c= n^2+ m^2= 2pq. Adding those equations, m^2+ 2mn+ n^2= (m+n)^2= p^2+ 2pq- q^2 so that p^2+2pq- q^2 is a perfect square. p^2+ 2pq- q^2= p^2+ 2pq+ q^2- 2q^2= (p+q)^2- 2q^2 so p and q must be such that (p+q)^2- 2q^2 is a perfect square.
So the whole question reduces to "can x^2- 2y^2 be a perfect square for any integers x and y?"
• Oct 24th 2009, 06:50 AM
aidan
Quote:
Originally Posted by Kobby
My first post and hope it is not a stupid one and related to this thread:
Can there be 4 whole numbers a, b, c, d all greater than zero such that:
a^2 + b^2 = c^2 and b^2 + c^2 = d^2?
If there are- can one give an example. If none can one prove it is not possible?
$
a^2 + b^2 = c^2$
$
b^2 + c^2 = d^2
$
replacing $c^2$ in second equation
$
b^2 + a^2 + b^2 = d^2
$
rearrange/simplify
$
a^2 + 2b^2 = d^2
$
$
995^2 + 2(100^2) = 1005^2
$
(There are others.)
• Oct 24th 2009, 08:47 AM
Kobby
Still 995, 100, 1005 are not pytagorial numbers- so still no solution.
The fact that there is a^2 + 2b^2 = d^2, doesnt give you a, b, c as pytagorial numbers.
What all you all written I figured already.
Are there any a, b, c, d out there?
• Oct 24th 2009, 11:46 AM
Opalg
Quote:
Originally Posted by Kobby
Can there be 4 whole numbers a, b, c, d all greater than zero such that:
a^2 + b^2 = c^2 and b^2 + c^2 = d^2?
If there are- can one give an example. If none can one prove it is not possible?
This is an interesting problem. I believe that these equations have no solutions. Suppose that there is a solution. Then after dividing through by any common factor we may assume that (a,b,c) and (b,c,d) are both primitive pythagorean triples. Thus c and d are odd, so b is even and a is odd. Then there are co-prime integers m (even) and n (odd) such that $b=2mn$ and $c=m^2+n^2$. Therefore $d^2 = b^2+c^2 = 4m^2n^2 + (m^2+n^2)^2 = m^4+6m^2n^2+n^2$.
I believe (I haven't checked it carefully) that the diophantine equation $m^4+6m^2n^2+n^2=d^2$ has no solutions in positive integers. The proof would be by Fermat's method of infinite descent (the same approach that Fermat himself used to show that $x^4+y^4=z^2$ has no solutions).
Write the equation as $(m^2+3n^2)^2 - d^2 = 8n^4$, and factorise the left side: $(m^2+3n^2+d)(m^2+3n^2-d) = 8n^4$. If a prime factor of the odd number n divides both of those factors then it also divides their sum and hence divides m. But m and n are co-prime, so that cannot happen. Therefore the factors $m^2+3n^2+d$ and $m^2+3n^2-d$ must be $2s^4$ and $4t^4$ (either in that order or in the opposite order), where s,t are co-prime and $st=n$.
If $m^2+3n^2+d = 2s^4$ and $m^2+3n^2-d = 4t^4$ then $m^2+3n^2 = s^4+2t^4$. But $n=st$, and so $m^2 = s^4 - 3s^2t^2 + 2t^4 = (s^2-t^2)(s^2-2t^2)$. But those two factors are again co-prime, and therefore they must both be squares, say $s^2-t^2 = p^2$ and $s^2-2t^2 = q^2$, where $m=pq$. Thus $q^2+t^2=p^2$ and $t^2+p^2=s^2$. So we have a solution to the same pair of equations as for a,b,c and d, but with smaller numbers.
The other possibility is that $m^2+3n^2+d = 4s^4$ and $m^2+3n^2-d = 2t^4$, but this leads to a similar conclusion.
This process can now be repeated, to find a sequence of smaller and smaller solutions. But you cannot have an infinite, strictly decreasing sequence of positive integers. So we have a contradiction, and the conclusion has to be that the equations have no solution.
• Oct 24th 2009, 01:45 PM
Kobby
Thanks really. In the wikipedia value you had mentioned they also refered to the subject: "There are no Pythagorean triplets in which the hypotenuse and one leg are the legs of another Pythagorean triple."
So a puzzle is solved.
Glad I came across this site.(Happy) | 2017-04-26T04:15:08 | {
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https://mathhelpboards.com/threads/problem-involving-matrix-multiplication-and-dot-product-in-one-proof.6463/ | # Problem involving matrix multiplication and dot product in one proof!
#### gucci
##### New member
The problem is:
Let A be a real m x n matrix and let x be in R^n and y be in R^m (n and m dimensional real vector spaces, respectively). Show that the dot product of Ax with y equals the dot product of x with A^Ty (A^T is the transpose of A).
The way I went about starting this problem is to use the definitions where the definition of the dot product on real numbers is: the sum with k from 1 to n of ak * bk and the definition of matrix multiplication for the entries of Ax would each be of the form: sum from k=1 to k=n of Aik * xk1.
Hopefully that was clear enough, but what I come up with when I plug the second definition into the first is one sum inside another and it seems like either I'm missing something that I can simplify or I went in the wrong direction! Does anyone have any suggestions for what I could try here? Any help is appreciated
#### Bacterius
##### Well-known member
MHB Math Helper
[JUSTIFY]Have you tried relabelling the indices of your nested sum? That and using the fact that $A^T$ is the transpose of $A$ should lead you to a proof that the two expressions are the same.
What I mean is in your nested sum, reverse the two indices $i$ and $j$ (or whatever their name is) and then state that $A_{ij} = A^T_{ji}$, and that should give you the desired result.[/JUSTIFY]
Last edited:
#### Chris L T521
##### Well-known member
Staff member
The problem is:
Let A be a real m x n matrix and let x be in R^n and y be in R^m (n and m dimensional real vector spaces, respectively). Show that the dot product of Ax with y equals the dot product of x with A^Ty (A^T is the transpose of A).
The way I went about starting this problem is to use the definitions where the definition of the dot product on real numbers is: the sum with k from 1 to n of ak * bk and the definition of matrix multiplication for the entries of Ax would each be of the form: sum from k=1 to k=n of Aik * xk1.
Hopefully that was clear enough, but what I come up with when I plug the second definition into the first is one sum inside another and it seems like either I'm missing something that I can simplify or I went in the wrong direction! Does anyone have any suggestions for what I could try here? Any help is appreciated
Here's an alternative to index labeling; it involves using matrix multiplication and noting that you can express the dot product as follows:
$(A\mathbf{x})\cdot \mathbf{y} = \mathbf{y}^TA\mathbf{x}.$
Likewise, $(A^T\mathbf{y})\cdot\mathbf{x} = \mathbf{x}^TA^T\mathbf{y}=(\mathbf{y}^TA\mathbf{x})^T$. Since $\mathbf{x}^TA^T\mathbf{y}$ and $\mathbf{y}^TA\mathbf{x}$ are $1\times 1$ matrices, we must have
$(\mathbf{x}^TA^T\mathbf{y})^T =\mathbf{x}^TA^T\mathbf{y}\quad\text{and} \quad (\mathbf{y}^TA\mathbf{x})^T = \mathbf{y}^TA\mathbf{x}$
and thus
$\mathbf{x}^TA^T\mathbf{y} = (\mathbf{y}^TA\mathbf{x})^T=\mathbf{y}^TA\mathbf{x}.$
Therefore, it follows that
$(A\mathbf{x})\cdot\mathbf{y} = \mathbf{y}^TA\mathbf{x} = (\mathbf{x}^T A^T\mathbf{y})^T = \mathbf{x}^TA^T\mathbf{y} = (A^T\mathbf{y})\cdot\mathbf{x}$
which is what you were after.
I hope this makes sense!
#### Deveno
##### Well-known member
MHB Math Scholar
I would simplify the presentation by noting that for any 1x1 matrix $M$, $M^T = M$
(there is only the single diagonal element).
Thus:
$Ax \cdot y = (y^T)(Ax) = [(y^T)(Ax)]^T = (Ax)^T(y^T)^T = (x^TA^T)y = x^T(A^Ty) = A^Ty \cdot x = x \cdot A^Ty$
(this only works with REAL inner-product spaces, by the way)
This is pretty much the same as what Chris L T521 posted, but has fewer steps.
Working with just elements, we have, for:
$A = (a_{ij}), x = (x_1,\dots,x_m), y = (y_1,\dots,y_n)$
$\displaystyle Ax \cdot y = \sum_{i = 1}^m \left(\sum_{j = 1}^n a_{ij}x_j\right) y_i$
$= (a_{11}x_1 + \cdots + a_{1n}x_n)y_1 + \cdots + (a_{m1}x_1 + \cdots + a_{mn}x_n)y_m$
$= (a_{11}y_1 + \cdots + a_{m1}y_m)x_1 + \cdots (a_{1n}y_1 + \cdots + a_{mn}y_m)x_n$
(make sure you understand that each term $a_{ij}x_jy_i = a_{ij}y_ix_j$ only occurs once in each sum, we're just grouping them differently...in the first sum we're matching the row entry index of $A$ with the index of $y$, in the second sum we're matching the column index of $A$ with the index of $x$, and the transpose just switches rows with columns).
$\displaystyle = \sum_{j = 1}^n \left(\sum_{i = 1}^m a_{ji}y_i \right)x_j$
$= A^Ty \cdot x = x \cdot A^Ty$
#### gucci
##### New member
Here's an alternative to index labeling; it involves using matrix multiplication and noting that you can express the dot product as follows:
$(A\mathbf{x})\cdot \mathbf{y} = \mathbf{y}^TA\mathbf{x}.$
So I get why everything after this step would follow, and thank you all very much for that. I am just missing why this is a true assumption :-/
#### Chris L T521
##### Well-known member
Staff member
So I get why everything after this step would follow, and thank you all very much for that. I am just missing why this is a true assumption :-/
Well if $\mathbf{u},\mathbb{v}\in\mathbb{R}^m$ where $\mathbf{u}=\begin{pmatrix} u_1 \\ u_2 \\ \vdots \\ u_m\end{pmatrix}$ and $\mathbf{v}=\begin{pmatrix} v_1 \\ v_2 \\ \vdots \\ v_m\end{pmatrix}$, then we know that $\mathbf{u}\cdot \mathbf{v} = u_1v_1 + u_2v_2 + \ldots + u_mv_m$
(note that as a matrix, a scalar quantity is a $1\times 1$ matrix).
However, $\mathbf{u}$ and $\mathbf{v}$ are $m\times 1$ matrices; thus, if we were to express the dot product in terms of matrix multiplication, we must have one $m\times 1$ and one $1\times m$ matrix. Hence, we need to take the transpose of one of the vectors to accomplish this. With that said, we can then say that
$\mathbf{u}\cdot\mathbf{v} = \mathbf{v}^T\mathbf{u} = \begin{pmatrix} v_1 & v_2 & \cdots & v_m \end{pmatrix} \begin{pmatrix} u_1 \\ u_2 \\ \vdots \\ u_m\end{pmatrix} = \begin{pmatrix}v_1u_1 + v_2u_2 + \ldots + v_mu_m\end{pmatrix} = \begin{pmatrix} u_1 & u_2 & \cdots & u_m\end{pmatrix} \begin{pmatrix} v_1 \\ v_2 \\ \vdots \\ v_m\end{pmatrix} = \mathbf{u}^T\mathbf{v} = \mathbf{v}\cdot\mathbf{u}.$
I hope this clarifies things!
#### Deveno
##### Well-known member
MHB Math Scholar
I'd like to point out that $B = A^T$ is the ONLY $n\times m$ matrix that makes:
$Ax \cdot y = x \cdot By$ true for ALL $x \in \Bbb R^n, y \in \Bbb R^m$.
For suppose $B$ is such a matrix. Since it holds for ALL $x,y$ it certainly must hold when:
$x = e_j = (0,\dots,1,\dots,0)$ (1 is in the $j$-th place) <--this vector is in $\Bbb R^n$
$y = e_i = (0,\dots,1,\dots,0)$ (1 is in the $i$-th place) <--this vector is in $\Bbb R^m$
Now $Ae_j = (a_{1j},\dots,a_{mj})$ (the $j$-th column of $A$) so
$Ae_j \cdot e_i = a_{ij}$ (the $i$-th entry of the $j$-th column of $A$).
By the same token, $Be_i = (b_{1i},\dots,b_{ni})$ and
$e_j \cdot Be_i = b_{ji}$.
Comparing the two (and using that $Ae_j \cdot e_i = e_j \cdot Be_i$), we see that: $b_{ji} = a_{ij}$, that is $B = A^T$.
So the equation $Ax \cdot y = x \cdot A^Ty$ can be used to DEFINE the transpose. This is actually useful later on when one is trying to define things in a "coordinates-free" way (Vector spaces don't come equipped with a "preferred" basis, we have to pick one, which is somewhat arbitrary. If we can prove things without using a particular chosen basis, our proof is "more general", which is typically seen as a GOOD thing in mathematics). | 2020-09-20T23:57:03 | {
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https://math.stackexchange.com/questions/2743992/is-there-a-pattern-in-the-sequence-l-1-l-2-l-3-ldots | Is there a pattern in the sequence $l_1,l_2,l_3,\ldots$?
Wilson's theorem asserts the following statement: $$(n-1)!\equiv -1\pmod n\Leftrightarrow n\text{ is prime.}\tag1$$ This means that \begin{align} n&{ \ \mid} \ \ (n-1)!+1 \\ \Leftrightarrow n&{ \ \mid} \ \ (n-1)!+1-n \\ &=(n-1)!-(n-1) \\ &= (n-1)\big((n-2)!-1\big) \\ \therefore n&{ \ \mid} \ \ (n-2)!-1 \\ \Leftrightarrow (n-2)!&\equiv 1\pmod n.\tag2\end{align} Or, in the congruence, I could add $n$ to the right hand side, and then divide both sides by $n-1$, yeilding the same result.
I then asked myself, why wouldn't Wilson's theorem assert that $$(n-2)!\equiv 1\pmod n\Leftrightarrow n\text{ is prime,}$$ as opposed to $(1)$? It is more useful because $(n-2)! < (n-1)!$, so we can test this theorem with more primes without the factorial becoming too large as fast in $(1)$, and it is still just as interesting. However, I went here and found that it really is not that necessary to simplify the theorem.
Now, I have been trying to find some similar result for $(n-3)!$, and I found that $$(n-3)!\equiv \frac{n-1}{2}\pmod n\Leftrightarrow n\text{ is prime > 2.}$$
Main Question: If for some $k\in\mathbb{Z^+}$ and remainder $l_k\in\mathbb{Z}$, $$(n-k)!\equiv l_k\pmod n,\tag{n is prime}$$ is there a pattern in the sequence $l_1, l_2, l_3,\ldots$?
Do there exist similar congruences for $(n-k)!$ such that $k > 3$?
This post was inspired by this post.
• Don't abuse notation. It is not true tbat $(n-1)!+1=(n-1)!-(n-1).$ – Thomas Andrews Apr 19 '18 at 2:35
• $(n-1)!\equiv (-1)\cdot (n-2)!\pmod{n}$ all the time, so there really isn't a cost to computing $(n-1)!.$ It's just a multiplication of $-1\pmod{n}.$ – Thomas Andrews Apr 19 '18 at 2:37
• @ThomasAndrews I know that. But it's just because of divisibility that we can cancel the $n$. If $n\mid a + n$ then of course $n\mid a$. That doesn't mean $a = a + n$. I know what u mean though – Feeds Apr 19 '18 at 4:45
• It is a matter of taste whether residue $1$ is more useful than residue $-1$. Calculating $(n-2)!$ modulo $n$ is as difficult as calculating $(n-1)!$ modulo $n$ , so we still cannot find large primes with this criterion. And even if we calculate $(n-100)!$ modulo $n$ for very large $n$, it won't make a significant diffference in complexity. – Peter Apr 19 '18 at 19:13
Note that since $$n$$ is prime, we can define the following: $$y \equiv \frac{a}{x} \pmod{m} \iff n \mid xy-a \space (x,y,a\neq 0)$$
We can use the same properties of congruence such as adding, subtracting, multiplying and dividing by relatively prime constants.
Now, to find $$(n-2)! \pmod{n} \space$$, we basically need to find $$\frac{1}{n-1} \pmod{n}$$. We can instead write $$n-1 \pmod{n}$$ as $$-1 \pmod{n}$$. Thus, we have $$\frac{1}{-1} \equiv -1 \pmod{n}$$ which gives you that $$(n-2)! \equiv 1 \pmod{n}$$ using Wilson's Theorem.
The same idea goes on for $$(n-k)!$$. We are to evaluate $$\frac{1}{-k} \pmod{n}$$. For $$k=1$$, we trivially get $$-1$$. For $$k=2$$, considering $$n$$ as an odd prime, we know that $$2 \mid n+1$$. Thus, $$\frac{n+1}{2}$$ is an integer from which we can derive $$(n-3)! \pmod{n}$$.
Unfortunately, this process can no longer continue. A prime can be any value modulo another odd prime less than it except $$0$$. Thus, we cannot detect any pattern.
• Thank you for your answer, but I had already figured this out. ‘Tis why I didn’t put a bounty on it, but I forgot to actually answer my own question.... however, your answer is nearly exactly the same as was my own, so I am gonna give you a tick. I’ll also give you a $+100$ bounty just because. Congratulations! $(+1)$ $\color{green}{\checkmark}$ – Feeds Sep 28 '18 at 5:02
• Thank you. I'm glad I could be helpful. – Haran Sep 28 '18 at 7:50
• It only took $5$ months, hahah :P (well, a "formal" answer, that is) – Feeds Sep 28 '18 at 7:51
• If you are interested in understanding factorials modulo primes, there is a question I have asked regarding them. Any progress or insight would be appreciated. The question is: math.stackexchange.com/questions/2651733/… – Haran Sep 28 '18 at 7:55
• Thanks. I will check it out :) – Feeds Sep 28 '18 at 7:55 | 2019-08-18T23:49:52 | {
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http://math.stackexchange.com/questions/374576/finding-the-smallest-integer-n-such-that-11-21-3-1-n%e2%89%a59/374584 | # Finding the smallest integer $n$ such that $1+1/2+1/3+…+1/n≥9$?
I am trying to find the smallest $n$ such that $1+1/2+1/3+....+1/n \geq 9$,
I wrote a program in C++ and found that the smallest $n$ is $4550$.
Is there any mathematical method to solve this inequality.
-
See this section of the "harmonic series" entry in Wikipedia. – amWhy Apr 27 '13 at 19:55
What you are looking for is $H_n = \sum_{1 \le k \le n} k^{-1}$. It is known that $H_n \approx \ln n + \gamma - \dfrac{1}{12 n^2} + \dfrac{1}{120 n^4} - \dfrac{1}{252 n^6} + O\left(\dfrac{1}{n^8}\right)$, where $\gamma \approx 0.5772156649$ is Euler's constant.
-
... $H_n$ is the sum you are computing. – meh Apr 27 '13 at 19:52
This nails it: $e^{9-0.5772} \approx 4549.62$. – lhf Apr 27 '13 at 21:17
@julien Right, copied the signs wrong :-( In any case, the approximation without the series already gives enough to solve your question. This is an asymptotic series the error is less than the first omited term. – vonbrand Apr 27 '13 at 21:49
@julien, note that $B_{2 n} \sim (-1)^n \left( \dfrac{n}{\pi e} \right)^{2 n}$, so the series does not converge. It can be shown that the error is less than the first term left out. – vonbrand Apr 28 '13 at 13:16
Right, thanks. I found that on wikipedia afterwards. Can it really be shown that the error is less than the first term left out at every order? Will Jagy, who knows this formula pretty well, says this is for sure only for very little orders. And yes, it can be shown then. And even more. We have inequalities. That is the point of my answer, inspired by Will Jagy. And this is why it is not sufficient to just give an asymptotic expansion. The more difficult part comes after that. – 1015 Apr 28 '13 at 13:32
We have $$\sum_{k=1}^N \dfrac1k \sim \log(N) + \gamma + \dfrac1{N} - \dfrac1{12N^2} + \mathcal{O}(1/N^3)$$ We want the above to be equal to $9$, i.e., $$9 = \log(N) + \gamma + \dfrac1{N} - \dfrac1{12N^2} + \underbrace{\mathcal{O}(1/N^3)}_{\text{can be bounded by }1/N^3}$$ We have $\log(N) + \gamma < 9 \implies N \leq 4550$. Also, $\log(N) + \gamma +\dfrac1{N} > 9 \implies N \geq 4550$. Hence, $$N = 4550$$
Here is the asymptotic from Euler–Maclaurin. We have $$\sum_{k=a}^b f(k) \sim \int_a^b f(x)dx + \dfrac{f(a) + f(b)}2 + \sum_{k=1}^{\infty} \dfrac{B_{2k}}{(2k)!}\left(f^{(2k-1)}(b)-f^{(2k-1)}(a)\right)$$ In our case, we have $f(x) = \dfrac1x$, $a=1$ and $b=N$. This gives us $$\sum_{k=1}^N \dfrac1k \sim \log(N) + \gamma + \dfrac{1}{2N} - \sum_{k=1}^{\infty} \dfrac{B_{2k}}{2k}\dfrac1{N^{2k}}$$
-
Very nice, thanks a lot. – kiranovalobas Apr 27 '13 at 19:54
What if $O(1/N^3)=10^{10^{10}}/N^3$? I know it is not... – 1015 Apr 27 '13 at 20:03
@julien Valid point. But as you said, it is not. – user17762 Apr 27 '13 at 20:12
@WillJagy I am aware of the asymptotic exapnsion of the harmonic series...I am afraid you don't see my point. The problem is to know whether the remainder is $\leq \frac{1}{N^k}$ or $\leq \frac{10^{10^{10^{10}}}}{N^k}$. In the latter case, how do you find the first $n$ that makes your sequence greater than $9$ just looking at the truncation? – 1015 Apr 27 '13 at 20:25
@WillJagy Do you really not see my point? I did not mean to bother anyone... – 1015 Apr 27 '13 at 20:52
The idea was given to me by Will Jagy. I am just writing it down because he suggested I do so.
Let $H_n:=\sum_{k=1}^n\frac{1}{k}$ the $n$-th harmonic number. It is known (see here or user17762's answer above) that $$H_n=\ln n+\gamma+\frac{1}{2n}+O\left( \frac{1}{n^2}\right)$$ where $\gamma$ denotes Euler-Mascheroni constant. It is not sufficient to know that to answer your question. You need inequalities, or a precise control of the error term in the expansion. Here are two inequalities which suffice.
Claim: we have $$S_n:=\ln n+\gamma <H_n< \ln n+\gamma+\frac{1}{2n}=:T_n$$ for all $n\geq 1$.
Application: with the lhs, a calculator, and $\gamma\simeq 0.5772156649$, we find $H_{4550}>9.00009817>9$. With the rhs, we get $H_{4549}<8.99998828<9$. So $n=4550$ is indeed the first such $n$.
Proof of the claim: first, let us show that $x_n:=H_n-S_n$ decreases. Indeed $$x_n-x_{n+1}=-\ln n+\ln(n+1)-\frac{1}{n+1}=-\ln\left(1-\frac{1}{n+1} \right)-\frac{1}{n+1}>0$$ where the last inequality follows for instance from the study of $f(x)=-\ln(1-x)-x$ on $(0,1)$. Since $\lim x_n=0$, it follows that $x_n>0$ for all $n\geq 1$. That is the lhs inequality.
Now let $y_n:=H_n-T_n$ and let us check that it increases. We have $$y_{n+1}-y_n=\frac{1}{2(n+1)}+\frac{1}{2n}-\ln\left(1+\frac{1}{n}\right)>0$$ where the inequality follows from the fact that the corresponding function of $x$, instead of $n$, is decreasing on $(0,+\infty)$ with limit $0$ at $+\infty$. So $y_n$ is increasing with limit $0$. This means that $y_n<0$ for all $n\geq 1$, which proves the rhs inequality. QED.
-
@WillJagy Here it is. Actually, I took a lazier path which suffices. Please let me know if there is anything wrong. Thanks a lot for showing me these inequalities. – 1015 Apr 27 '13 at 23:42
That's right. The fact that this is decisive for threshold 9 is somewhat a matter of luck. The usual practice is to get the size of the window an extra order smaller then the last term added, that being $1/n$. In this case, that would mean considering your $T_n$ and $U_n = T_n - \frac{1}{12 n^2}.$ For example, what you have written does not suffice to find the first $H_n$ that exceeds 2, or exceeds 3, or 5, or 12. – Will Jagy Apr 28 '13 at 1:37
@WillJagy Yes. I started doing it with the order $2$ on the left and $4$ on the right. But I took a little break in doing so, which was more cumbersome, to try some computations. And I was surprised to see that these easier estimates suffice. – 1015 Apr 28 '13 at 2:14
Well, if you do one more step, (and I recommend just the one additional) you will find that, as is common for strictly alternating asymptotic series, the true value is bounded between consecutive terms in the asymptotic series. This is not guaranteed when all you have is an asymptotic series, but is guaranteed for strictly alternating convergent infinite series. So you see, I was never worried about the issue that bothered you. – Will Jagy Apr 28 '13 at 2:46
@WillJagy As I don't know Bernoulli numbers very well, I don't know that $\frac{|B_{2k}|}{2kN^{2k}}$ is nonincreasing. Otherwise, I would not have worried about that either as I know Leibniz criterion and its consequences. Is it really straightforward to see that the latter is nonincreasing? – 1015 Apr 28 '13 at 2:50
In case you do not understand vonbrand and user17762's reasoning, both are using the fact that $\int_1^n \frac{dk}{k}=\ln{n}$, and that a sum can be approximated by an integral (for large n values) with a correcting factor. Here the correcting factor is $\gamma$, the Euler-Mascheroni constant.
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You can look up "digamma function" for more info. A version I made up myself, which uses a symmetry property to give every other power, is $$H_n = \log \left( n + \frac{1}{2} \right) + \gamma + \frac{1}{6(2n+1)^2} - \frac{7}{60(2n+1)^4} + \frac{31}{126(2n+1)^6} - \frac{127}{120(2n+1)^8} + \frac{511}{66(2n+1)^{10}} - O \left( \frac{1}{(2n+1)^{12}} \right)$$ with $$\gamma = 0.5772156649015328606065...$$
I put part of this in my programmable calculator, just the 1/6 blah - 7/60 blahblah. Because of the alternating signs, this is definitive, too small with $n=4549,$ big enough with $n=4550.$ There is no guarantee the $\pm$ signs continue to alternate, and the rational coefficients are allowed to grow, note $511/66 \approx 7.74.$ The use of asymptotic expansions is to truncate at a convenient place and know enough about the possible error. Interesting to put this on a high precision program, with extremely high precision demanded for $\gamma,$ and look at the error of what I wrote multiplied by $(2n+1)^{12}.$
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@julien, I think you should post your own answer. – Will Jagy Apr 27 '13 at 20:47 | 2016-04-29T19:50:19 | {
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https://www.physicsforums.com/threads/average-chord-length-of-a-circle.929479/ | # I Average chord length of a circle
1. Oct 24, 2017
### serverxeon
I would like to find the average chord length of a circle.
And I have 2 methods, which gave different answers...
[The chord is defined as the line joining 2 points on the circumference of the circle.]
The general formula for a chord length is $d=2R\sin(\delta/2)=2\sqrt{R^2-u^2}$
Method 1: Integrate over angles
Avg length = $$\frac{\int_0^{2\pi} 2R\sin(\delta/2)\,d\delta}{2\pi}$$
$$=\frac{4R}{\pi}$$
Method 2: Integrate over diameter (-R to R)
Avg length = $$\frac{\int_{-R}^{R} 2\sqrt{R^2-u^2}\,du}{2R}$$
After simplification
$$\int_{-R}^{R} \sqrt{1-\left(\frac{u}{R}\right)^2}\,du$$
Then, using the fact
$$\int_{-Z}^{Z} \sqrt{1-\left(\frac{u}{Z}\right)^2}\,du = \frac{Z\pi}{2}$$
I get
$$=\frac{R\pi}{2}$$
So the answers from my 2 methods dont add up. Any things i might have overlooked?
2. Oct 24, 2017
### Staff: Mentor
You found two different definitions for the average. The question "what do you average over?" is not just interesting for the calculation, it leads to different mathematical questions - with different answers.
3. Oct 24, 2017
### serverxeon
I am not too sure. Are they not the same?
However, even if I ignore the different definition of average and solely look at the 'sum of all chords',
I get (for angles) $8R$ and (for diamater) $R^2\pi$
Am I not simply summing up all possible chord lengths from 0, to 2R, then back to 0 again?
4. Oct 24, 2017
### Staff: Mentor
There is an infinite set of chords, there is no unique way to “sum” over them.
An analogy: Let $a=b^2$. Clearly a=b has the solutions a=b=0 and a=b=1. But the integrals between these limits are different: $$\int_0^1 a da \neq \int_0^1 b^2 db$$
The key point here is $da \neq db$ - it matters how you integrate.
5. Oct 24, 2017
### serverxeon
Your example is not exactly correct.
If you substituted $$a=b^2$$
then
$$\int_0^1 a\, da = \int_0^1 {b^2\cdot2b}\,db$$
which is then the same...
Anyway,
Did you mean to say there is ONE set of infinite chords, or there are infinite sets of finite chords?
Because you seem to suggest the former, and that would mean there should only be one unique answer?
6. Oct 24, 2017
### mathman
There is one set of infinite chords. However, there are many possible ways to define the probability distribution. Your analysis uses two different ones, so it is not surprising to get two different answers for the mean.
7. Oct 25, 2017
### Staff: Mentor
Yes, but you had to introduce the factor 2 b. That's the point. A change in coordinates can change the integral if you do not add factors like these.
8. Nov 2, 2017
### disregardthat
The one true way to obtain a random chord, is to uniformly pick two points P,Q on the circle. These choices amounts to a uniform choice of angles $\theta,\phi \in [0,2\pi)$. The angle POQ, where O is the center, is equal to $\sqrt{2R^2-2R^2\cos(\theta-\phi)} = \sqrt{2}R\sqrt{1-\cos(\theta - \phi)}$ by the cosine formula (note that the sign of $\theta - \phi$ does not matter here). The average length thus becomes
$$\frac{\int^{2\pi}_0 \int^{2\pi}_0 \sqrt{2}R\sqrt{1-\cos(\theta - \phi)}d\theta d \phi}{\int^{2\pi}_0 \int^{2\pi}_0 d\theta d\phi} = \frac{\int^{2\pi}_0 \int^{2\pi}_0 \sqrt{2}R\sqrt{1-\cos(\theta - \phi)}d\theta d \phi}{(2\pi)^2}$$
Evaluating the inner integral first, it is clear that as $\phi$ varies as a constant, the integral remains unaffected. So
$$\int^{2\pi}_0\left( \int^{2\pi}_0 \sqrt{2}R\sqrt{1-\cos(\theta - \phi)}d\theta \right) d \phi = \int^{2\pi}_0\left( \int^{2\pi}_0 \sqrt{2}R\sqrt{1-\cos(\theta)}d\theta \right) d \phi$$
$$= 2\pi \sqrt{2}R \int^{2\pi}_0 \sqrt{1-\cos(\theta)}d\theta = 2\pi \sqrt{2}R \left( \int^{\pi}_0 \sqrt{1-\cos(\theta)}d\theta + \int^{\pi}_0 \sqrt{1+\cos(\theta)}d\theta \right)$$
On the last equality we split the integral in two, and used the fact that $\cos(\theta + \pi) = -\cos(\theta)$. Next observe that $1-\cos(\theta) = 2\sin^2(\frac{\theta}{2})$ and $1 + \cos(\theta) = 2\sin^2(\frac{\theta}{2})$, so that $\sqrt{1-\cos(\theta)} = \sqrt{2}\sin(\frac{\theta}{2})$ and $\sqrt{1+\cos(\theta)} = \sqrt{2}\cos(\frac{\theta}{2})$ (note that both expressions to the right are positive since we now are integrating over $\theta \in [0,\pi]$!).
Our integral becomes
$$2\pi \sqrt{2}R \int^{\pi}_0 \sqrt{2}\left( \sin(\frac{\theta}{2}) + \cos(\frac{\theta}{2}) \right) d\theta = 4\pi R \int^{\pi}_0 \left( \sin(\frac{\theta}{2}) + \cos(\frac{\theta}{2}) \right)d\theta$$
This is easily evaluated, yielding $4\pi R[-2\cos(\frac{\theta}{2})+2\sin(\frac{\theta}{2})]^{\pi}_0 = 16 \pi R$. By our first formula, we obtain the average length
$$\frac{16 \pi R}{(2\pi)^2} = \frac{4R}{\pi}$$ which indeed was your first answer.
Last edited: Nov 2, 2017 | 2018-03-17T17:12:15 | {
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http://am207.info/wiki/montecarlointegrals.html | Keywords: monte-carlo | integration | uniform distribution | law of large numbers | lotus | central limit theorem | normal distribution | Download Notebook
Contents
%matplotlib inline
import numpy as np
import matplotlib.pyplot as plt
import seaborn as sns
The basic idea
Let us formalize the basic idea behind Monte Carlo Integration in 1-D.
Consider the definite integral:
Consider:
If $V$ is the support of the uniform distribution on a to b then the pdf $U_{ab}(x) = \frac{1}{V} = \frac{1}{b-a}$
Then from LOTUS and the law of large numbers:
or
Practically speaking, our estimate will only be as exact as the number of samples we draw, but more on this soon..
Example.
Calculate the integral $I= \int_{2}^{3} [x^2 + 4 \, x \,\sin(x)] \, dx.$
We know from calculus that the anti-derivative is $x^3/3 + 4\sin(x) -4x\cos(x).$
To solve this using MC, we simply draw $N$ random numbers from 2 to 3 and then take the average of all the values $f(x)=x^2 + 4 \, x \,\sin(x)$ and normalized over the volume; this case the volume is 1 (3-2=1).
def f(x):
return x**2 + 4*x*np.sin(x)
def intf(x):
return x**3/3.0+4.0*np.sin(x) - 4.0*x*np.cos(x)
a = 2;
b = 3;
# use N draws
N= 10000
X = np.random.uniform(low=a, high=b, size=N) # N values uniformly drawn from a to b
Y =f(X) # CALCULATE THE f(x)
V = b-a
Imc= V * np.sum(Y)/ N;
exactval=intf(b)-intf(a)
print("Monte Carlo estimation=",Imc, "Exact number=", intf(b)-intf(a))
Monte Carlo estimation= 11.8120823531 Exact number= 11.8113589251
Mutlidimensional integral.
That is nice but how about a multidimensional case?
Let us calculate the two dimensional integral $I=\int \int f(x, y) dx dy$ where $f(x,y) = x^2 +y^2$ over the region defined by the condition $x^2 +y^2 ≤ 1$
In other words we are talking about a uniform distribution on the unit circle
fmd = lambda x,y: x*x + y*y
# use N draws
N= 8000
X= np.random.uniform(low=-1, high=1, size=N)
Y= np.random.uniform(low=-1, high=1, size=N)
Z=fmd(X, Y) # CALCULATE THE f(x)
R = X**2 + Y**2
V = np.pi*1.0*1.0
N = np.sum(R<1)
sumsamples = np.sum(Z[R<1])
print("I=",V*sumsamples/N, "actual", np.pi/2.0) #actual value (change to polar to calculate)
I= 1.56308724855 actual 1.5707963267948966
Monte-Carlo as a function of number of samples
How does the accuracy depends on the number of points(samples)? Lets try the same 1-D integral $I= \int_{2}^{3} [x^2 + 4 \, x \,\sin(x)] \, dx$ as a function of the number of points.
Imc=np.zeros(1000)
Na = np.linspace(0,1000,1000)
exactval= intf(b)-intf(a)
for N in np.arange(0,1000):
X = np.random.uniform(low=a, high=b, size=N) # N values uniformly drawn from a to b
Y =f(X) # CALCULATE THE f(x)
Imc[N]= (b-a) * np.sum(Y)/ N;
plt.plot(Na[10:],np.sqrt((Imc[10:]-exactval)**2), alpha=0.7)
plt.plot(Na[10:], 1/np.sqrt(Na[10:]), 'r')
plt.xlabel("N")
plt.ylabel("sqrt((Imc-ExactValue)$^2$)")
#
<matplotlib.text.Text at 0x1140aa780>
Obviously this depends on the number of $N$ as $1/\sqrt{N}$.
Errors in MC
Monte Carlo methods yield approximate answers whose accuracy depends on the number of draws. So far, we have used our knowledge of the exact value to determine that the error in the Monte Carlo method approaches zero as approximately $1/\sqrt{N}$ for large $N$, where $N$ is the number of trials.
But in the usual case, the exact answer is unknown. Why do this otherwise?
So, lets repeat the same evaluation $m$ times and check the variance of the estimate.
# multiple MC estimations
m=1000
N=10000
Imc=np.zeros(m)
for i in np.arange(m):
X = np.random.uniform(low=a, high=b, size=N) # N values uniformly drawn from a to b
Y =f(X) # CALCULATE THE f(x)
Imc[i]= (b-a) * np.sum(Y)/ N;
plt.hist(Imc, bins=30)
plt.xlabel("Imc")
print(np.mean(Imc), np.std(Imc))
11.8114651823 0.00398497853806
This looks like our telltale Normal distribution.
This is not surprising
Estimating the error in MC integration using the CLT.
We know from the CLT that if $x_1,x_2,…,x_n$ be a sequence of independent, identically-distributed (IID) random variables from a random variable $X$, and that if $X$ has the finite mean $\mu$ AND finite variance $\sigma^2$.
Then,
converges to a Gaussian Random Variable with mean $\mu$ and variance $\sigma^2/n$ as $n \to \infty$:
This is true regardless of the shape of $X$, which could be binomial, poisson, or any other distribution.
The sums
are exactly what we want to calculate for Monte-Carlo Integration(due to the LOTUS) and correspond to the random variable f(X) where X is uniformly distributed on the support.
Whatever the original variance of f(X) might be, we can see that the variance of the sampling distribution of the mean goes down as $1/n$ and thus the standard error goes down as $1/\sqrt{n}$ as we discovered when we compared it to the exact value as well.
Why is this important?
Comparing to standard integration techniques
What if we changed the dimensionality of the integral? The formula for $S_n$ does not change, we just replace $g(x_i)$ by $g(x_i, y_i, z_i…)$. Thus, the CLT still holds and the error still scales as $\frac{1}{\sqrt{n}}$.
On the other hand, if we divide the $a, b$-interval into $N$ steps and use some regular integration routine, what is the error? Consider the midpoint rule as illustrated in this diagram from Wikipedia:
The basic idea is that the function value at the midpoint of the interval is used as the height of the approximating rectangle. In general, the differing methods consist of choosing different $x_i$ below..with left being at the left end, right being at the right end. $I(est) = \sum_i f(x_i)\Delta x_i = \frac{b-a}{n} \sum_i f(x_i)$
The error on the estimation of the integral can be shown to decrease as $\frac{1}{n^2}$. The basic reason for this can be understood on a taylor series expansion of the function to second order. When you integrate on the sub-interval, the linear term vanishes while the quadratic term becomes cubic in $\Delta x$. So the local error goes as $\frac{1}{n^3}$ and thus the global as $\frac{1}{n^2}$.
Monte-Carlo if clearly not competitive with the midpoint method in 1-D. Its actually not even competitive with left or right rectangle methods.
The trapeziod rule uses a line between the sub-interval points while the Simpsons rule uses a quadratic.
These integrations can be generalized to multiple dimensions, and the rule for these
• left or right rule: $\propto 1/n$
• Midpoint rule: $\propto 1/n^2$
• Trapezoid: $\propto 1/n^2$
• Simpson: $\propto 1/n^4$
where $n=N^{1/d}$. MC becomes better than the Simpson method only in 8 dimensions.. | 2019-01-17T22:55:32 | {
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http://mathhelpforum.com/pre-calculus/72958-distance-time-word-problem.html | # Math Help - distance and time word problem
1. ## distance and time word problem
A man walks for 45 minutes at a rate of 3 mph, then jogs for 75 minutes at a rate of 5 mph, then sits and rests for 30 minutes, and finally walks for 1 hour and half.Find the rule of the function that expresses this distance traveled as a function of time [ Caution: Don't mix up the units of time ; use either minutes or hours , not both]
could someone help me with this problem? I sketched the graph ,and the ranges I got are from 0 to .75 , .75 to 2, 2 to 2.5 and 2.5 to 4 . I know that the first function is 3t and the third one is 6.25.However, I don't know how to write the other two ones.
2. Originally Posted by vance
A man walks for 45 minutes at a rate of 3 mph, then jogs for 75 minutes at a rate of 5 mph, then sits and rests for 30 minutes, and finally walks for 1 hour and half.Find the rule of the function that expresses this distance traveled as a function of time [ Caution: Don't mix up the units of time ; use either minutes or hours , not both]
could someone help me with this problem? I sketched the graph ,and the ranges I got are from 0 to .75 , .75 to 2, 2 to 2.5 and 2.5 to 4 . I know that the first rule is 3t and the third one is 6.25.However, I don't know how to write the other two ones.
Hello!
The function rule for this problem is
distance = rate x time
d = rt
To write the rule as the distance traveled as a function of the time it took to travel that distance you will need to evaluate each set of data using this function. Since you have 3 different rates of speed, and 4 time periods, you will have to calculate each one.
I changed the times to hours.
RATE TIME
0 $\rightarrow$ .5
3 $\rightarrow$ .75
5 $\rightarrow$ 1.25
3 $\rightarrow$ 1.5
So, distance as a function of time:
d= rt
1) $d(t) = 0t$ so $d(.5) = (0)(.5) = 0 miles$
2) $d(t) = 3t$ so $d(.75) = (3)(.75) = 2.25 miles$
3) $d(t) = 5(t)$ so $d(1.25) = (5)(1.25) = 6.25 miles$
4) $d(t) = 3t$ so $d(1.5) = (3)(1.5) = 4.5 miles$
Hope this is what you needed! Good luck!
3. Hello, vance!
This is a tricky one!
Luckily, I've done one of these before . . .
A man walks for 45 minutes at a rate of 3 mph, then jogs for 75 minutes at 5 mph,
then sits and rests for 30 minutes, and finally walks for 1 hour and half.
Find the rule of the function that expresses this distance traveled as a function of time.
This is a piece-wise function.
The distance function depends on what time period is involved.
From $t=0\text{ to }t = 0.75$, the distance function is: . $3t$
. . By the end of the 45 minutes, he has gone 2.25 miles.
From $t = 0.75\text{ to } t = 2$, the distance function is: . $2.25 + 5t$
. . By the end of the 75 minutes, he has gone another 6.25 miles, a total of 8.5 miles.
From $t = 2\text{ to }t = 2.5$, the distance function is: . $8.5 + 0t \:=\:8.5$
. . By the end of the 30 minutes, he hasn't moved; his total is still 8.5 miles.
From $t = 2.5\text{ to }4$, the distance function is: . $8.5+3t$
. . By the end of the 90 minutes, he has gone another 4.5 miles, a total of 13 miles.
The function would be written like this:
. . $d(t) \;=\;\begin{Bmatrix}3t & & 0 \leq t \leq 0.75 \\ 2.25 + 5t & & 0.75 \leq t < 2 \\ 8.5 & & 2 \leq t \leq 2.5 \\ 8.5+3t & & 2.5 < t < 4 \end{Bmatrix}$
The graph would look like this:
Code:
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4. Originally Posted by Soroban
Hello, vance!
This is a tricky one!
Luckily, I've done one of these before . . .
This is a piece-wise function.
The distance function depends on what time period is involved.
From $t=0\text{ to }t = 0.75$, the distance function is: . $3t$
. . By the end of the 45 minutes, he has gone 2.25 miles.
From $t = 0.75\text{ to } t = 2$, the distance function is: . $2.25 + 5t$
. . By the end of the 75 minutes, he has gone another 6.25 miles, a total of 8.5 miles.
From $t = 2\text{ to }t = 2.5$, the distance function is: . $8.5 + 0t \:=\:8.5$
. . By the end of the 30 minutes, he hasn't moved; his total is still 8.5 miles.
From $t = 2.5\text{ to }4$, the distance function is: . $8.5+3t$
. . By the end of the 90 minutes, he has gone another 4.5 miles, a total of 13 miles.
The function would be written like this:
. . $d(t) \;=\;\begin{Bmatrix}3t & & 0 \leq t \leq 0.75 \\ 2.25 + 5t & & 0.75 \leq t < 2 \\ 8.5 & & 2 \leq t \leq 2.5 \\ 8.5+3t & & 2.5 < t < 4 \end{Bmatrix}$
The graph would look like this:
Code:
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Thanks soroban. This is what I was looking for!!! | 2015-03-27T00:13:17 | {
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http://math.stackexchange.com/questions/851830/is-it-okay-to-reverse-engineer-proofs-in-homework-questions | # Is it okay to reverse engineer proofs in homework questions?
In a linear algebra text book, one homework question I received was:
Prove that $\mathbf{a \cdot b} = \frac{1}{4}(\|\mathbf{a + b}\|^2 - \|\mathbf{a - b}\|^2)$.
Where $\mathbf{a}$ and $\mathbf{b}$ are vectors in $\Bbb{R}^n$.
This is trivial to prove if we start from $\frac{1}{4}(\|\mathbf{a + b}\|^2 - \|\mathbf{a - b}\|^2)$ and reverse engineer it in $\Bbb{R}^2$: $$\|\mathbf{a + b}\|^2 = a_1^2 + 2a_1b_1 + b_1^2 + a_2^2 + 2a_2b_2 + b_2^2 \\ \|\mathbf{a - b}\|^2 = a_1^2 - 2a_1b_1 + b_1^2 + a_2^2 - 2a_2b_2 + b_2^2 \\ \|\mathbf{a + b}\|^2 - \|\mathbf{a - b}\|^2 = 4a_1b_1 + 4a_2b_2 \\ \frac{1}{4}(\|\mathbf{a + b}\|^2 - \|\mathbf{a - b}\|^2) = \frac{4}{4}(a_1b_1 + a_2b_2) \\ = a_1b_1 + a_2b_2 = \mathbf{a \cdot b}$$
But I'm worried about whether or not proofs like this are "legal", if that makes any sense. There was no wording in the question stating that I couldn't start from the right side of the identity, but I still have this strange feeling of guilt that I should've tried solving the identity starting from the left side and working in the "normal" direction.
For questions like these, is it okay to start from the right side of the identity? Would what I get out of doing the question in reverse be the same as if I did it normally?
-
Yes, of course that's legal. In math, you're allowed to do anything that's logically justified. – anomaly Jun 30 '14 at 2:43
In proving identities, it's perfectly acceptable to start from the RHS and go to the LHS, as long as your steps are always fully "reversible". That means that you can use the double-implication symbol $\iff$ between every step. So you might have to be careful about doing stuff like going from $x^2 = a^2 \implies x = a$, because that's not acceptable ($x$ can be $-a$) but if you've been careful to avoid that sort of pitfall, it's fine. However, if you're worried about the presentation, you should recognise that you should be able to reverse everything and rewrite it to go from the LHS to RHS. – Deepak Jun 30 '14 at 2:49
I've done stuff like this many times - as I'm sure many others have around these parts. Sometimes it's hard to see how one side relates to the other but maybe you can make it meet in the middle. – Cameron Williams Jun 30 '14 at 2:51
@Deepak It seems like you are confusing logical equivalence with equality. The OP is asking about a proof that two things are equal. – Trevor Wilson Jun 30 '14 at 5:39
OP: What Trevor and David mean is that, when you prove $a = b$ by finding intermediate $c$s, you don't have to worry about "reversible steps", because if $b = c$, then $c = b$. But if you want to prove $a = b$, and you want to "reverse engineer" it by messing with both sides of the equation, and deriving a truth, then you worry about reversibility, because $P \implies Q$ isn't the same as $Q \implies P$. – Henry Swanson Jun 30 '14 at 14:05
It's absolutely fine to reverse engineer proofs! As long as the proof works in the end, it doesn't matter how you got to it - in fact, reverse engineering proofs is a fairly standard technique.
However you need to be careful in this approach: you need to make sure that each step is reversible ($p$ implies $q$ does not mean $q$ implies $p$). With your example here though, we can easily reverse each step because everything is just equality.
It is worth noting that although your proof method is fine, the proof itself doesn't quite get what you want because you have assumed that $a,b\in\mathbb{R^2}$ when you wish to prove for any two vectors $a,b\in\mathbb{R^n}$. This problem, though, is easily dealt with (can you see how?).
-
Equality is symmetric, by definition. As such, any concerns one might have about reversibility apply just as much when an equality is proven "left-to-right" as when an equality is proven "right-to-left". It's a little misleading to suggest that these concerns only come up when going in one direction. – David Richerby Jun 30 '14 at 9:39
I read the "be careful" as applying to the more general question about "proofs like this," which might or might not be applied to propositions of the form $A=B$. "With your example here though ... everything is just equality." In other words, the fact that the proposition is of the form $A=B$ addresses the concern about reversibility. – David K Jun 30 '14 at 11:56
Somewhat relevant xkcd comic on how you can take it a bit too far: Handy exam trick: when you know the answer but not the correct derivation, derive blindly forward from the givens and backward from the answer, and join the chains once the equations start looking similar. Sometimes the graders don't notice the seam. ;) On a serious note though, working from both ends is an extremely useful approach. – Voo Jun 30 '14 at 23:26
Personally, I'd suggest this approach of solution (which address the loss of generality by assuming $\mathbb{R^2}$:
$$4\mathbf{a}\cdot \mathbf{b} = [\mathbf{(a+b)\cdot(a+b)} - \mathbf{(a-b)\cdot(a-b)}]$$
which leads almost immediately to the required result.
-
Not only is it 'OK', it is equally as valid to go from the RHS to the LHS as it is to go from the LHS to the RHS.
Moreover, sometimes one finds that the best approach is to work from both ends simultaneously.
As a student you may well find that the simultaneous approach is very useful in a exam situation.
Perhaps when fully written up, a single direction might be easier to follow, but the overwhelming concern is correctness.
Indeed, one thing that might impress a maths examiner beyond correctness is cleverness, which in maths usually means finding a trick or shortcut. From that viewpoint, being smart about thinking whether to from RHS to LHS rather than LHS to RHS is a good thing.
-
I don't know about "reverse engineering" in general, but I'd guess it's always fine: as anomaly commented, you can do anything that's logically justified.
In this case the answer is particularly clear, however. Equality is a symmetric relation: $x = y$ means exactly the same thing as $y = x$. A proof of one is the same as a proof of the other (with the addition of a step at the end where you use this symmetry property to flip the two sides, but usually there is no reason to write this because it will be obvious to the reader.)
I can imagine some situations where writing $x = y$ rather than $y = x$ might be better from a stylistic perspective, but style is usually less important that correctness when doing homework.
-
It is definitely a valid proof, as the other answers mention.
One downside with a proof like this is that you are left wondering how you come up with the identity in the first place..
- | 2015-10-07T12:42:54 | {
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https://math.stackexchange.com/questions/938065/t-be-linear-operator-on-v-by-tb-ab-ba-prove-that-if-a-is-a-nilpotent-ma | # $T$ be linear operator on $V$ by $T(B)=AB-BA$. Prove that if A is a nilpotent matrix, then $T$ is a nilpotent operator.
Let $V$ be a vector space of $n\times n$ matrices over a field F, and let $A$ be a fixed $n\times n$ matrix. $T$ be linear operator on $V$ by $T(B)=AB-BA$. Prove that if A is a nilpotent matrix, then $T$ is a nilpotent operator.
I have done in a way that $T^2(B)= T(AB-BA)= A^2B-2ABA+BA^2\Rightarrow T^3(B)=A^3B-3A^2BA+3ABA^2-BA^3$ Proceeding in this way since A is nilpotent $\exists m\in\Bbb N$ s.t $T^m(B)=0$.
Hence T is also a nilpotent operator. This a problem of Hoffman Kunze of Primary Decomposition Chapter. So can anyone give me any other solution because this solution depends on some sort of intution. Please don't use any Rational and Jordan Forms formula.
• Ok but how? Calculating the matrix is much laborious. I have done it for $2\times2$ case. If u have any idea u can write in the answer. – user152715 Sep 19 '14 at 17:18
• No, my question is different. – user152715 Sep 19 '14 at 17:28
• No, its not. There may be some steps in between this u can see but I can't. If u derive the matrix of T for $2\times2$ case. U will see the matrix of T does directly depend upon the elements of A but not directly as the structure of the matrix A. So the eigenvalue they have derived can't be derived in that way. Anyway if u still thinks that there is a relation then post 1 answer deriving the eigenvalue. – user152715 Sep 19 '14 at 17:41
• I'm getting confused between my As and Bs. – copper.hat Sep 19 '14 at 17:53
• Give the answer reading the question – user152715 Sep 19 '14 at 18:32
This is in the same spirit to your proof, but presented in a different way. If $\lambda B = AB-BA$ for some $B\ne0$ and some $\lambda$ in the algebraic closure of $F$, then $(A-\lambda I)B=BA$ and $(A-\lambda I)^k B=BA^k$ for any $k\ge1$. In particular, $(A-\lambda I)^nB=0$. However, if $\lambda$ is nonzero, $A-\lambda I$ would be invertible and hence $B=0$, which is a contradiction.
• Can U explain the line from $(A-\lambda I)B=BA$ to $(A-\lambda I)^kB=BA^k$. – user152715 Sep 20 '14 at 3:45
• @user152715 Try to prove by mathematical induction. E.g. $(A-\lambda I)^2B=(A-\lambda I)\left[(A-\lambda I)B\right]=(A-\lambda I)BA=\left[(A-\lambda I)B\right]A=BA^2$. – user1551 Sep 20 '14 at 8:38
• Why $A-\lambda I$ would be invertible for nonzero $\lambda$? – Majid May 20 '18 at 20:19
• @Majid As $A$ is nilpotent, all eigenvalues of $A-\lambda I$ are equal to $-\lambda\ne0$. Hence $A$ is invertible. Alternatively, as $A$ is nilpotent, $A^m=0$ for some $m\ge1$. So, if $(A-\lambda I)x=0$, then $0=A^mx=\lambda^m x$ and hence $x=0$. Therefore $\ker A=0$ and $A$ is invertible. – user1551 May 20 '18 at 20:27
• @user1551 Thanks! I got. – Majid May 20 '18 at 20:29
This is not an answer as it depends on certain characteristics of the underlying field (also, I haven't used the fact that $A$ is nilpotent).
Suppose $T(B) = \mu B$, that is $\mu B = AB-BA$. Then $\mu B^2 = BAB-B^2A= (AB-\mu B)B-B^2 A= AB^2 -B^2A - \mu B^2$, or $2 \mu B^2 = AB^2 -B^2 A = T(B^2)$.
So, if $\mu$ is an eigenvalue corresponding to an eigenvector $B$ of $T$, then $B^2$ is an eigenvector corresponding to the eigenvalue $2 \mu$. Hence $2^k \mu$ are eigenvalues.
Since there are only a finite number of eigenvalues (this is where I am making presumptions about the field), we have $\mu = 0$.
Your proof seems correct to me. This result is used for Engel's theorem in the theory of Lie algebras. If $x\in \mathbb{gl}(V)$ is nilpotent, then also $ad(x)$ is nilpotent, where $ad(x)(y)=[x,y]=xy-yx$ for $x,y\in \mathfrak{gl}(V)$. Indeed, $ad(x)^m$ is a linear combination of terms $x^iyx^{m-i}$.
I proved this the same way the OP proved this problem (coincidentally, I had the same problem as homework, proved it, and then went online to see if another proof was available, to verify my work -- as I worked all the way up to the expression [I'm guessing the OP did this as well]). $T^{m}(B)={\displaystyle{\sum\limits_{i=0}^{m}(-1)^{i}\binom{m}{i}A^{m-i}BA^{i}}}$, then factor $B$ out of the sum, which gives the result, for any $B\in\mathcal{M}_{n\times n}(F)$ -- etc. Further, I did happen to find another proof as well. I'll give a proof sketch, since, to me, this alternative proof is convoluted a little, or somewhat overkill; mine/OP's proof is shorter, less complicated, etc. Nevertheless, here is the proof-sketch.
Proof-Sketch: First, if $V$ is an $F$-vector space, we claim that if $T_{1},T_{2}\in\mathcal{L}(V,V)$ are nilpotent, commuting operators, then $T_{1}+T_{2}$ is nilpotent.
To prove the claim, we have $T_{1}^{k}=0_{\mathcal{L}(V,V)}=T_{2}^{m}$ for some $k,m\in\mathbb{N}$ since these operators are assumed nilpotent (suppose without loss of generality that $m\leq k$ and note that $0_{\mathcal{L}(V,V)}$ is the zero operator/element of $\mathcal{L}(V,V)$ which is my preference in notation). Use the Binomial Theorem to expand $(T_{1}+T_{2})^{k+m}$, and show that $T_{1}^{k+m-i}T_{2}^{i}=0_{\mathcal{L}(V,V)}$ for all $i=0,1,...,m,m+1,...,k+m$; note that the use of the Binomial Theorem in this case is possible since the operators commute.
Returning to the overall proof, we have $V=\mathcal{M}_{n\times n}(F)$ whenever $n\in\mathbb{N}$ is arbitrarily fixed, and we take $A\in\mathcal{M}_{n\times n}(F)$ arbitrarily also. Define the operators $T_{1},T_{2}\in\mathcal{L}(V,V)$ such that for all $B\in\mathcal{M}_{n\times n}(F)$ we have $T_{1}(B)=AB$ and $T_{2}(B)=-BA$. This all being said, we first need to show that on $\mathcal{M}_{n\times n}(F)$ we have:
(i) $T_{1}T_{2}=T_{2}T_{1}$ (i.e., $T_{1},T_{2}$ commute); and
(ii) $T_{1}+T_{2}=T$ (where $T$ is defined the OP's problem-statement - also, this follows trivially by the definition of a sum of operators).
Thus, to show $T$ is nilpotent, it is sufficient to show that $T_{1},T_{2}$ are both nilpotent operators on $\mathcal{M}_{n\times n}(F)$. This can be accomplished by claiming for any $r,s\in\mathbb{N}$, we have $T_{1}^{r}(B)=A^{r}B$ as well as $T_{2}^{s}(B)=(-1)^{s}BA^{s}$ and proceeding by induction on $r$ and $s$ in $\mathbb{N}$, respectively. Thus, since $T_{1},T_{2}$ are nilpotent, we apply the claim at the beginning of the proof to get the overall result, and we are done.
$~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~\square$
Like I said, to me this seems to be a bit much, since we have to show a bit more that any of the other proofs posted, but I decided to post it since the original claim at the beginning of the proof can be useful (note that claim is independent of the dimension of $V$). Hopefully there aren't any mistakes, as I tried to keep the above proof-sketch a little abstract so one who decides to use it can work out the details himself/herself (and I was a little quick about things too...sorry about that and that this is a bit lengthy still yet).
• To critique my own work, where I say $T^{m}(B)={\displaystyle{\sum\limits_{i=0}^{m}\binom{m}{i}A^{m-i}BA^{i}}}$, I think I need $T^{m}(B)={\displaystyle{\sum\limits_{i=0}^{m}(-1)^{i}\binom{m}{i}A^{m-i}BA^{i}}}$ instead (as it is mentioned above, $T$ is defined in the OP's problem-statement)? Going to edit this in now, awaiting any corrections, thoughts, etc. – Procore Jul 12 '18 at 23:55
Note that $$T=L-R \qquad L=\textrm{Multiply X at Left}\qquad R=\textrm{Multiply X at right}$$ Since $$L\circ R=R\circ L\qquad R^{n}=0\qquad L^{n}=0$$ if $A^n=0$, then $T^{2n}=0$ by binomial theorem. | 2019-06-19T01:40:18 | {
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https://math.stackexchange.com/questions/2842530/showing-if-fracxc-sim-textgammaa-b-then-x-sim-textgammaa-cb | # Showing if $\frac{X}{c}\sim\text{Gamma}(a,b)$ then $X\sim\text{Gamma}(a,cb)$
I am trying to show that if $$\frac{X}{c}\sim\text{Gamma}(a,b) \ \ \ \ \ \ \text{then} \ \ \ \ \ \ \ X\sim\text{Gamma}(a,cb)$$
My first approach was to use a PDF transformation. I let $$Z=\frac{X}{c}\Rightarrow X=Zc$$ Then $$f_X(x)=f_Z\Big(\frac{x}{c}\Big)\Big|\frac{dz}{dx}\Big|=\frac{1}{\Gamma(a)b^a}e^{-\frac{x}{cb}}\Big(\frac{x}{c}\Big)^{a-1}\frac{1}{c}=\frac{1}{\Gamma(a)(cb)^a}e^{-\frac{x}{cb}}x^{a-1}$$ Which is clearly the density function of $\text{Gamma}(a,cb)$. Hence $X\sim\text{Gamma}(a,cb)$.
But, I know tried to confirm this using an MGF approach.
$$m_z(u)=\mathbb{E}\Big(e^{Zu}\Big)=\mathbb{E}\Big(e^{\frac{X}{c}u}\Big)=m_X\Big(\frac{u}{c}\Big)=\Big(1-b\frac{u}{c}\Big)^{-a}$$ which is the MGF for the $\text{Gamma}\Big(a,\frac{b}{c}\Big)$ distribution.
Where have I made my mistake in the MGF approach?
• $m_Z(u)=m_X(u/c)$. – StubbornAtom Jul 6 '18 at 5:55
• Typo fixed. Thanks for spotting that! My result still follows though – user557493 Jul 6 '18 at 5:56
• Which gamma pdf are you using? – StubbornAtom Jul 6 '18 at 5:59
• Oh I see, I assumed that $X\sim\text{Gamma}(a,b)$ which is incorrect. What's the best way to show $X\sim\text{Gamma}(a,cb)$ using MGFs? – user557493 Jul 6 '18 at 6:03
• If $X\sim\text{Gamma}(a,b)$, then $$f(x)=\frac{1}{\Gamma(a)b^a}e^{-\frac{x}{b}}x^{a-1} \ \ \ x>0$$ – user557493 Jul 6 '18 at 6:36
Let $$Z = \dfrac{X}{c} \sim \text{Gamma}(a, b)$$.
Then $$M_{Z}(u) = (1-bu)^{-a} = \mathbb{E}[e^{uZ}]\tag{*}$$ and thus, observing that $$X = cZ$$,
$$M_{X}(u)=\mathbb{E}[e^{uX}] = \mathbb{E}[e^{ucZ}]=\mathbb{E}[e^{(uc)Z}]=M_{Z}(uc) = [1-b(uc)]^{-a}=[1-(bc)u]^{-a}$$ hence $$X \sim \text{Gamma}(a, bc)$$.
Where did you go wrong in your work? First of all,
$$m_z(u) \neq \Big(1-b\frac{u}{c}\Big)^{-a}$$ We can see that since $$Z \sim \text{Gamma}(a, b)$$ that its MGF should be given by (*) above.
Second of all, $$m_X\Big(\frac{u}{c}\Big) \neq \Big(1-b\frac{u}{c}\Big)^{-a}$$ What you did here - assuming you were doing the work from left to right - was you assumed that $$X \sim \text{Gamma}(a, b)$$ to begin with (how else would you know what $$m_{X}(u)$$ is equal to that?), but we actually don't have that assumption available to us. That is, you shouldn't make any assumptions about what $$M_{X}$$ is. You know what $$M_{Z}$$ is because you have an assumption of what the distribution of $$Z$$ is available to you, but this is not the case for $$X$$.
As usual, the calculations depend on whether the distribution is parametrized by shape and rate; i.e., $$f_X(x) = \frac{b^a x^{a-1} e^{-bx}}{\Gamma(a)}, \quad M_X(t) = (1 - t/b)^{-a},$$ or shape and scale: $$f_X(x) = \frac{x^{a-1} e^{-x/\theta}}{\theta^a \Gamma(a)}, \quad M_X(t) = (1 - \theta t)^{-a}.$$ If you use the shape-rate parametrization, then $Y = X/c$ has density $$f_Y(y) = c f_X(cy),$$ which implies $Y$ has rate $cb$. This is corroborated by the MGF approach: $$M_Y(t) = \operatorname{E}[e^{tY}] = \operatorname{E}[e^{(t/c)X}] = M_X(t/c) = (1 - t/(cb))^{-a}.$$ If you use the shape-scale parametrization, you get $$f_Y(y) = cf_X(cy)$$ as before, but this means $Y$ has scale $\theta/c$. The MGF is $$M_Y(t) = M_X(t/c) = (1 - (\theta/c) t)^{-a},$$ which is again consistent.
It is worth noting that I have done the reverse transformation; that is to say, you assumed $X/c \sim \operatorname{Gamma}(a,b)$ and I have taken $X \sim \operatorname{Gamma}(a,b)$. My $Y$ is your $X$, and my $X$ is your $Z$.
• So, are you saying that if $$\frac{Y^2}{\sigma^2}\sim\text{Gamma}(a,b)\Rightarrow Y^2\sim\text{Gamma}\Big(a,b\sigma^2\Big)??$$ – user557493 Jul 6 '18 at 9:01 | 2019-08-24T18:47:46 | {
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https://math.stackexchange.com/questions/2740468/deriving-properties-of-the-exp-function-from-definition | # Deriving properties of the exp function from definition
Define exponential function $\exp$ as follows: $$\exp:\Bbb{R}\rightarrow(0,\infty)$$ \begin{align}i) \ \forall x,y\in \Bbb{R}:\exp(x+y)=\exp(x)\cdot\exp(y)\end{align} \begin{align}ii) \ \forall x\in \Bbb{R}:\exp(x)\ge x+1 \end{align} Book says these 4 properties can be derived using this definition.
$$(1) \ \exp \text{is increasing on } \Bbb{R}$$ $$(2) \ \exp(0)=1$$ $$(3) \ \forall x\in \Bbb{R}: \forall n \in \Bbb{Z}: \exp(nx)=(\exp(x))^n$$ $$(4) \ \forall x\in \Bbb{R},x>0:\exp(x)>1$$ The only thing i managed to do myself was to prove $(3)$ by mathematical induction for $n\in\Bbb{N}$. I also tried to prove the monotonicity directly from definition, but can't seem to find a starting point. Would you please give me hint on where to start with proofs of these? Another thing I am interested in is a proof of uniqueness of this definition. How do i prove, assuming only these two definitions $i),ii)$ that $\exp$ function defined like this is unique?
EDIT: For the proof of $(2)$ i have following idea: $$\exp(0)=\exp{(1+(-1))}=\exp(1)\cdot\exp(-1)$$ and by $(3)$ (not proven yet for negative integers) we can continue $$\exp(1)\cdot\exp(-1)=\frac{\exp(1)}{\exp(1)}=1$$ EDIT 2: What remains is to show that $(3)$ holds for negative values of $n$.
• From (i), (1) and (4) are equivalent. – Lord Shark the Unknown Apr 16 '18 at 21:19
By your condition (ii), $\exp(x) > 1$ for all $x > 0$. Hence if $y > x$, we have
$$\exp(y) = \exp(x + (y - x)) = \exp(x) \cdot \exp(y - x) > \exp(x) \cdot 1$$
giving monotonicity. The fact that $\exp(0) = 1$ follows from setting $x = y = 0$ and knowing that $\exp(0) > 0$.
As for uniqueness, the standard technique is to take two candidates and compare them. For example, if $\exp_2$ is another function satisfying the conditions, study $\exp - \exp_2$ and make sure it's the zero function. Or something more natural here is to take $\exp / \exp_2$ and show it's identically one.
• @MichalDvořák I fixed a typo in my answer. It should have been about $x = y = 0$. – user296602 Apr 16 '18 at 21:31
• I would be interested in an elementary argument showing that it is indeed fixed by the two conditions above. I really do think you need at least some real analysis here. – Thomas Bakx Apr 16 '18 at 21:57
I have never come across this definition of the exponential. Where did you find this? Apart from that, let me help you:
Note that we can take $x=y=0$ in i) to find $\mbox{exp}(0)=1$ (it cannot be zero because of ii) as you can see). Also, if we have (4) then we also have (1) by using (i) again. In fact, $\mbox{exp}$ is then strictly increasing. So it remains to prove (4).
If there is an $x>0$ such that $\mbox{exp}(x)=1$, then by (3), which we already proved, we obtain a set of arbitrarily large real numbers $nx$ for which $\mbox{exp}(nx)=1$. This must contradict ii).
As for uniqueness, I think you need something a little more sophisticated. Can you prove that $\mbox{exp}$ is differentiable with derivative equal to itself, using the properties you established? If this is the case, you know that its derivative at $0$ must equal one because otherwise property (ii) cannot hold. Observe that this property ii) is actually very restrictive, because for example if we only use (1)-(4) as assumptions, any function of the form $x \mapsto a^x$ would fit the bill for positive $a$.
Edit: as an elaboration, if we know that the function is differentiable with derivative equal to itself, then the quotient of any two functions satisfying i) and ii) must have zero derivative by application of the quotient rule, so that it is identically $1$.
• I've found that in scripts from Faculty of Mathematics and Physics from the Charles university of Prague, unfortunately, not available in English version – Michal Dvořák Apr 16 '18 at 22:03
• If i can prove that $\exp$ is its own derivative, isn't only thing i could show by the quotient rule that, assuming we have two different $\exp$'s $$\bigg(\frac{\exp_1(x)}{\exp_2(x)}\bigg)'=0$$ Can i actually assume that $\frac{\exp_1(x)}{\exp_2(x)}=1$? What makes sense to me is to assume that $\frac{\exp_1(x)}{\exp_2(x)}=c$ for some $c\in\Bbb{R}$ Also unfortunately, i can't seem to prove the derivative, i still need to know somehow shot, that $$\lim_{h\to 0}\frac{\exp(h)-1}{h}=1$$, which i have no idea how to show withouht knowing the limit definition for e and logarithms. – Michal Dvořák Apr 16 '18 at 22:20
• As for the quotient rule: true, but both functions must satisfy i) and ii), so both also satisfy (2). This should tell you what the value of $c$ is. For the derivative, you indeed need to calculate the limit you mentioned. Note that by using ii) twice, you can show that $\mbox{exp}(h) \geq 1+h$ and $\mbox{exp}(-h) \geq 1-h$, or by (3) this is $\mbox{exp}(h) \leq \frac{1}{1-h}$. Now apply the squeeze theorem for limits. – Thomas Bakx Apr 16 '18 at 22:22
• Oh, I'm sorry. You should apply i) to find $\mbox{exp}(h)\mbox{exp}(-h)=1$, which yields $\mbox{exp}(-h) = \frac{1}{\mbox{exp}(h)}$. Also I should have been a little more careful: the inequality $\mbox{exp}(h) \leq \frac{1}{1-h}$ only holds for $h<1$ of course, but we are only interested in values of $h$ close to zero anyways. – Thomas Bakx Apr 16 '18 at 22:33
• This is not so hard to see: we have, in your notation, that $\mbox{exp}_1(0) = \mbox{exp}_2(0) = 1$. We saw that the quotient is constant, so it equals its value at $0$. It follows that $c = \frac{1}{1}=1$. – Thomas Bakx Apr 16 '18 at 22:43 | 2019-06-19T00:45:04 | {
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http://mathhelpforum.com/trigonometry/164691-larger-print.html | Which is larger?
• November 28th 2010, 07:19 PM
Mr Rayon
Which is larger?
Let x be an element of R. Find which term has a larger value, sin(cos x) or cos(sin x)?
Any help or detailed working out would be appreciated!
• November 30th 2010, 05:29 AM
Sudharaka
Quote:
Originally Posted by Mr Rayon
Let x be an element of R. Find which term has a larger value, sin(cos x) or cos(sin x)?
Any help or detailed working out would be appreciated!
Dear Mr Rayon,
Please refer the attached file. The blue line represents y=sin(cosx) while the green curve represents y=cos(sinx). From this graph it is clear that sin(cosx)<cos(sinx). But for the moment I dont have any idea of how prove this result.
• November 30th 2010, 07:04 AM
Quote:
Originally Posted by Mr Rayon
Let x be an element of R. Find which term has a larger value, sin(cos x) or cos(sin x)?
Any help or detailed working out would be appreciated!
Here is another view.
$sinx$ ranges from $-1\rightarrow\ 1$
$cos(sinx)$ therefore ranges from $cos(-1)=cos(1)=0.54\rightarrow\ cos(0)=1$
$cosx$ ranges from $-1\rightarrow\ 1$
$sin(cosx)$ therefore ranges from $sin(-1)=-0.84\rightarrow\ sin(1)=0.84$
A little more analysis will show they are never equal.
• November 30th 2010, 03:41 PM
Sudharaka
Quote:
Here is another view.
$sinx$ ranges from $-1\rightarrow\ 1$
$cos(sinx)$ therefore ranges from $cos(-1)=cos(1)=0.54\rightarrow\ cos(0)=1$
$cosx$ ranges from $-1\rightarrow\ 1$
$sin(cosx)$ therefore ranges from $sin(-1)=-0.84\rightarrow\ sin(1)=0.84$
A little more analysis will show they are never equal.
Just out of curiosity, with what software did you create your images???
• November 30th 2010, 03:49 PM
Quote:
Originally Posted by Sudharaka
Just out of curiosity, with what software did you create your images???
Hi Sudharaka,
I only use the "pages" application in Apple iWorks!
in conjunction with a Mathtype program (very basic).
AM
• December 1st 2010, 12:42 AM
BobP
Quote:
Originally Posted by Mr Rayon
Let x be an element of R. Find which term has a larger value, sin(cos x) or cos(sin x)?
Any help or detailed working out would be appreciated!
Both functions are continuous and you can't find a value of x for which they are equal.
So, if one is bigger than the other for a particular value of x, it will be bigger for all values of x.
• December 1st 2010, 04:39 AM
Here is a way we can develop the proof...
$cos(sinx)>0$
hence we need not be concerned with $sin(cosx)<0$
and the maximum value of $sin(cosx)$ is <1.
We only need consider both functions >0 which is in quadrant 1, where $0
In quadrant 1, $\displaystyle\ sin(cosx)=cos\left(\frac{\pi}{2}-cosx\right)$
Hence we are examining $\displaystyle\ cos(sinx)-cos\left(\frac{\pi}{2}-cosx\right)$
so we want to know if $sinx<\left[\frac{\pi}{2}-cosx\right]$
since in quadrant 1, $cos(angle)>cos(larger\;angle)$
$\frac{\pi}{2}>sinx+cosx$ ?
Hence we find the maximum value of $sinx+cosx$ by differentiating...
$cosx-sinx=0\Rightarrow\ cosx=sinx$ which happens at $x=\frac{\pi}{4}$ radians.
The maximum value of $sinx+cosx=\frac{1}{\sqrt{2}}+\frac{1}{\sqrt{2}}=\f rac{2}{\sqrt{2}}=\sqrt{2}$
$\frac{\pi}{2}>\sqrt{2}\Rightarrow\ cos(sinx)>sin(cosx)$
• December 2nd 2010, 04:49 AM
Sudharaka
Quote:
Originally Posted by BobP
Both functions are continuous and you can't find a value of x for which they are equal.
So, if one is bigger than the other for a particular value of x, it will be bigger for all values of x.
Hi BobP,
Is this a theorem? Can you please tell me where you have encountered it?
• December 4th 2010, 12:52 AM
BobP
Quote:
Originally Posted by Sudharaka
Hi BobP,
Is this a theorem? Can you please tell me where you have encountered it?
If it isn't a theorem I'm happy to claim it !
Actually, I thought it was just common sense (Happy).
• December 4th 2010, 06:38 AM
Sudharaka
Hi everyone,
Below is a proof that I have formulated for the theorem BobP had given in post #6. If you find any mistake in it please do not hesitate to tell me. Using this result the above problem could be easily solved.
• December 4th 2010, 07:14 AM
The graph of the two functions $sin(cosx)$ and $cos(sinx)$ or $cos(sinx)-sin(cosx)$ are a nice guide.
If we are presented with just the 2 functions, $cos(sinx)$ and $sin(cosx)$
and we want to compare them, how do we answer "what stands out ?"
• December 4th 2010, 04:40 PM
Sudharaka
Quote:
The graph of the two functions $sin(cosx)$ and $cos(sinx)$ or $cos(sinx)-sin(cosx)$ are a nice guide.
If we are presented with just the 2 functions, $cos(sinx)$ and $sin(cosx)$
and we want to compare them, how do we answer "what stands out ?"
I dont quite understand what you meant by saying ".......how do we answer what stands out?" Can you please explain this.
• December 4th 2010, 07:57 PM
SammyS
Which is larger?
Quote:
Originally Posted by BobP
Both functions are continuous and you can't find a value of x for which they are equal.
So, if one is bigger than the other for a particular value of x, it will be bigger for all values of x.
This tread is getting old, but Mr. Ryan still seemed confused.
This may be a Lemma to BobP's theorem.
I hope you can agree that: If two functions, f(x) and g(x) are continuous, and if they are not equal for any x, then one of them is greater than the other for all x.
In this case let's see if we can solve:
$\sin(\cos(x))=\cos(\sin(x))$.
$\cos(\theta)=\sin(\theta+{\pi\over2})$,
so $\sin(\cos(x))=\sin(\sin(x)+{\pi\over2})$.
Then $\cos(x)=\sin(x)+{\pi\over2}$. (O.K., you could add integer multiple of $2\pi$ to this, but that won't help.)
$\cos(x)-\sin(x)={\pi\over2}$
Use the identity: $\cos(x)-\sin(x)=\sqrt{2}\cos(x+{\pi\over4})$.
This gives us $\sqrt{2}\cos(x+{\pi\over4})={\pi\over2}$
The maximum value of the cosine is 1, and $\sqrt{2}<{\pi\over2}$, so $\sin(\cos(x))=\cos(\sin(x))$ has no solution.
Since $\sin(\cos(\pi))=\sin(-1)\approx -0.8415$ (It's definitely negative!)
and $\cos(\sin(\pi))=\cos(0)=1$,
it must be that $\sin(\cos(x))<\cos(\sin(x))$, for all x.
• December 31st 2010, 05:52 PM
Maybe simplest and without reference to graphs is:
$cos(A+B)=cosAcosB-sinAsinB\Rightarrow\ cos\left(\frac{\pi}{2}-A\right)=sinA$
$cos(sinx)-sin(cosx)=cos(sinx)-cos\left(\frac{\pi}{2}-cosx\right)$
This is the difference of cosines.
$cosA-cosB=-2sin\left[\frac{A+B}{2}\right]sin\left[\frac{A-B}{2}\right]$
$\Rightarrow\ cos(sinx)-sin(cosx)=-2sin\left[\frac{sinx-cosx+\frac{\pi}{2}}{2}\right]sin\left[\frac{sinx+cosx-\frac{\pi}{2}}{2}\right]$
One could go further with trigonometry here.
An alternative is to examine the contents of the square brackets on the previous line.
We can examine the maximum and minimum values of these expressions.
$\frac{d}{dx}(sinx-cosx)=cosx+sinx=0\Rightarrow\ cosx=-sinx$
Referring to the unit-circle, this occurs when $x=\frac{3{\pi}}{4},\;\;x=\frac{7{\pi}}{4}$
(1) maximum value within the first square brackets
$x=\frac{3{\pi}}{4}\Rightarrow\frac{\frac{1}{\sqrt{ 2}}+\frac{1}{\sqrt{2}}+\frac{\pi}{2}}{2}=\frac{1}{ \sqrt{2}}+\frac{\pi}{4}$
$0<\frac{1}{\sqrt{2}}+\frac{\pi}{4}<\frac{\pi}{2}\R ightarrow\ sin\left[\frac{sinx-cosx+\frac{\pi}{2}}{2}\right]>0$
(2) minimum value within the first square brackets
$x=\frac{7{\pi}}{4}\Rightarrow\frac{-\frac{1}{\sqrt{2}}-\frac{1}{\sqrt{2}}+\frac{\pi}{2}}{2}=-\frac{1}{\sqrt{2}}+\frac{\pi}{4}$
$0<-\frac{1}{\sqrt{2}}+\frac{\pi}{4}<\frac{\pi}{2}\Rig htarrow\ sin\left[\frac{sinx-cosx+\frac{\pi}{2}}{2}\right]>0$
$\frac{d}{dx}(sinx+cosx)=cosx-sinx=0\Rightarrow\ cosx=sinx\Rightarrow\ x=\frac{\pi}{4},\;\;x=\frac{5{\pi}}{4}$
(3) maximum value within the second square brackets
$x=\frac{\pi}{4}\Rightarrow\frac{\frac{1}{\sqrt{2}} +\frac{1}{\sqrt{2}}-\frac{\pi}{2}}{2}=\frac{1}{\sqrt{2}}-\frac{\pi}{4}$
$-\frac{\pi}{2}<\frac{1}{\sqrt{2}}-\frac{\pi}{4}<0\Rightarrow\ sin\left[\frac{sinx+cosx-\frac{\pi}{2}}{2}\right]<0$
(4) minimum value within the second square brackets
$x=\frac{5\pi}{4}\Rightarrow\frac{-\frac{1}{\sqrt{2}}-\frac{1}{\sqrt{2}}-\frac{\pi}{2}}{2}=-\frac{1}{\sqrt{2}}-\frac{\pi}{4}$
$-\frac{\pi}{2}<-\frac{1}{\sqrt{2}}-\frac{\pi}{4}<0\Rightarrow\ sin\left[\frac{cosx+sinx-\frac{\pi}{2}}{2}\right]<0$
Therefore, it follows that
$cos(sinx)-sin(cosx)>0\Rightarrow\ cos(sinx)>sin(cosx)$ | 2016-09-01T04:40:20 | {
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https://math.stackexchange.com/questions/2115380/please-critique-my-proof-that-sqrt12-is-irrational | Please critique my proof that $\sqrt{12}$ is irrational
I would like critiques on correctness, conciseness, and clarity. Thanks!
Proposition: There is no rational number whose square is 12
Proof: Suppose there were such a number, $a = \in \mathbb{Q}$ s.t. $a^2 = 12$.
This implies that $\exists$ $m, n \in \mathbb{Z}$ s.t. $\frac{m^2}{n^2} = 12.$ Assume without loss of generality that $m,~ n$ have no factors in common.
$\Rightarrow m^2 = 12n^2$.
This implies that $m^2$ is even, and therefore that $m$ is even; it can thus be written $2k = m$ for some $k \in \mathbb{Z}$.
Thus $m^2 = 12n^2$
$\Rightarrow 4k^2 = 12n^2$
$\Rightarrow \frac{k^2}{3} = n^2$
Because $n^2$ is an integer, it is clear that $3$ divides $k^2$ which imples that $k$ has $3$ or $\frac{k}{n}$ has a factor (because $\frac{k^2}{n^2}= 3$)
Suppose that the former is true, and $3$ is a factor of $k$. Then $k = 3j$ for some integer j, which implies that $(3j)^2 = 3n^2$
$\Rightarrow 9j^2 = 3n^2$
$\Rightarrow n^2 = 3j^2$
$\Rightarrow n^2 = \frac{k^2}{n^2}j^2$
$\Rightarrow k = \frac{n^2}{j}$ but this implies that $j$ divides $n^2$, but $j$ divides $m$, and by initial assumption $n$ and $m$ have no factors in common, so this is a contradiction.
Suppose now that $\frac{k}{n}$ is a factor of k. Then $k = \frac{k}{n}j$ for some integer $j$. Then $(\frac{k}{n}j)^2 = 3n^2$ which implies that $3j^2 = 3n^2 \Rightarrow j^2 = n^2 \Rightarrow j = n$. But this means that $n$ divides $m$, which again is a contradiction. Thus any rational representation of the number whose square equals $12$ leads to a contradiction and this number must therefore have no rational representation.
• Looks fine to me. A shorter version: $\sqrt{12}=2\sqrt{3}$ belongs to $\mathbb{Q}$ iff $\sqrt{3}$ belongs to $\mathbb{Q}$, but that is impossible by the unique factorization theorem and the primality of $3$. – Jack D'Aurizio Jan 26 '17 at 18:45
• Correct thinking, but you can get away with a lot less work and writing.:) To begin with, write: $$\sqrt{12} = 2 \sqrt{3}.$$ This means, all you need is to show that $\sqrt{3}$ is irrational. – avs Jan 26 '17 at 18:46
• Looks good. An alternate approach is to look at the polynomial $x^2-12$ and use the rational root theorem to show that any square root of $12$ would have to be a divisor of $12$. – lulu Jan 26 '17 at 18:47
• 3 divides k^2 which imples that k has 3 or k/n has a factor I don't understand the or part. Since $3 \mid k^2$ it follows that $3 \mid k\,$ and there is no or case. – dxiv Jan 26 '17 at 18:52
• @BenL For any integers $a=bc \implies b \mid a$ and $c \mid a\,$. That's by definition, and requires no additional assumptions. You have $k^2 = 3 n^2\,$, therefore $3 \mid k^2\,$. – dxiv Jan 26 '17 at 18:59
Proof. Assume $\sqrt{12} \in \mathbb{Q}$ is rational, then it can be written as $\sqrt{12}=\cfrac{m}{n}$ with $m,n \in \mathbb{Z}$ coprime.
Squaring the equality gives $m^2 = 12 n^2 = 3 \cdot 4 \cdot n^2\,$. Therefore $3 \mid m^2 = m \cdot m$ and, since $3$ is a prime, it follows by Euclid's Lemma that $3 \mid m\,$.
Then $m = 3k$ for some $k \in \mathbb{Z}$ and substituting back gives $9 k^2 = 12 n^2 \iff 3 k^2 = 4 n^2\,$. Therefore $3 \mid 4 n^2$ and, since $3 \not \mid 4$ it follows that $3 \mid n^2$ then, again by Euclid's Lemma, $3 \mid n\,$.
But $3 \mid m$ and $3 \mid n$ contradicts the assumption that $m,n$ are coprime, so the premise that $\sqrt{12} \in \mathbb{Q}$ must be false, therefore $\sqrt{12}$ is irrational.
Critique of the posted proof.
Proof: Suppose there were such a number, $a = \in \mathbb{Q}$ s.t. $a^2 = 12$.
This implies that $\exists$ $m, n \in \mathbb{Z}$ s.t. $\frac{m^2}{n^2} = 12.$ Assume without loss of generality that $m,~ n$ have no factors in common.
$\Rightarrow m^2 = 12n^2$.
So far so good.
This implies that $m^2$ is even, and therefore that $m$ is even;
The fact that $2 \mid m^2 \implies 2 \mid m$ may sound obvious, but still needs some justification. You could argue by contradiction, or use Euclid's Lemma.
it can thus be written $2k = m$ for some $k \in \mathbb{Z}$.
Thus $m^2 = 12n^2$
$\Rightarrow 4k^2 = 12n^2$
Correct. As an observation, $k^2 = 3 n^2$ just eliminated the perfect square factor of $4$ and reduced the problem to proving that $\sqrt{3}$ is irrational.
$\Rightarrow \frac{k^2}{3} = n^2$
Because $n^2$ is an integer, it is clear that $3$ divides $k^2$ which imples that $k$ has $3$
You should generally avoid fractions where they are not necessary. The previous line gave $k^2 = 3 n^2\,$, which directly implies that $3 \mid k^2\,$.
or $\frac{k}{n}$ has a factor (because $\frac{k^2}{n^2}= 3$)
This makes no sense, and it is in fact not needed to complete the proof.
Suppose that the former is true, and $3$ is a factor of $k$. Then $k = 3j$ for some integer j, which implies that $(3j)^2 = 3n^2$
$\Rightarrow 9j^2 = 3n^2$
$\Rightarrow n^2 = 3j^2$
The proof is complete right here at this point, if you just note that the last equality implies that $3 \mid n^2\,$, and therefore $3 \mid n$ which contradicts the assumption that $m,n$ are coprime.
[ rest of post snipped ]
• My main concern is proving that $3 | n^2 \Rightarrow 3|n$. I'm familiar with Euclid's Lemma, but I'm not clear whether it's legit in the context of a real analysis course. Thanks though, this was very helpful! – BenL Jan 26 '17 at 20:01
• @BenL Euclid's Lemma is of course true, regardless of context. Here, $n$ is an integer, and $n^2=n\cdot n$ is a product of two integers (which happen to be equal). $3$ is a prime which divides that product, so it must divide (at least) one of the factors. Ergo, $3 \mid n\,$. – dxiv Jan 26 '17 at 20:03
• @BenL You don't need Euclid's Lemma for that. mod $3\!:\ n\not\equiv 0\,\Rightarrow\,n\equiv \pm1\,\Rightarrow\, n^2\equiv 1\,\Rightarrow\,3\nmid n^2.\$ But you do need Euclid (or equivalent) if you want to prove that for *all* primes $p$ since then you cannot brute force check all possible residues as above. – Bill Dubuque Jan 26 '17 at 20:03
"This implies that ∃ m,n∈Z s.t. m2n2=12. Assume without loss of generality that m, n have no factors in common."
I, personally, would not argue "without loss of generality" . $q \in \mathbb Q$ is defined as $q = \frac mn$ for some relatively prime integers. So we declare them to have no factors in common by fiat-- not merely by lack of loss of generality. (It's not that big of an issue.)
"This implies that $m^2$ is even, and therefore that m is even"
I'd accept this but dxiv very much has a point, that it should require some justification. I personally would simply put it in more definitive language. I'd say: "Therefore $2|m^2$ and, as $2$ is prime, $2|m$". This could require a little justification in that all numbers have a unique prime factorization so that for prime $p$ we know if $p|ab$ then $p|a$ or $p|b$ so if $p|m^2$ then $p|m$ or $p|m$.
"Because $n^2$ is an integer, it is clear that 3 divides $k^2$ which implies that k has 3 or $\frac kn$ has a factor ".
As $\frac {k^2}3$ is a integer, it implies $3|k^2$. Period. That always happens. That any thing else may happen doesn't matter. It may have $k/n$ as a factor or it may have $7$ as a factor. Or it may not. Those don't matter.
Also, if $\frac kn$ is an integer at all, then it is trivial that $\frac kn$ is a factor $k$ whether or not $k^2/3$ is an integer or not. And if $\frac kn$ is not an integer then the statement $\frac kn$ is a factor of $k$ is meaningless.
And if $k/n$ is an integer, then $n|m = 2k$ and as $n,m$ have no factor in common then $n = 1$. (Which would mean $\sqrt{12} = 2\sqrt{3}$ is an integer which is easy to verify is not the case).
"if the former ($3|k$) then .... "
All that is just fine and the rest is unneeded.
But the rest is a bit of a mess.
"Suppose now that $k/n$ is a factor of $k$" Again, this is trivial if $n|k$ and is meaningless if $n \not \mid k$.
And we can rule out $n|k$ as that would imply $\sqrt{12} = m/n = 2k/n$ is an integer. Which
But beware. This is true of all numbers and nothing relevant is likely to arise. And it doesn't:
" Then $(\frac knj)^2=3n^2$ which implies that $3j^2=3n^2$"
Actually, no, it implies $3j^2 = 3n^4$. And thus we get $j = n^2$ (we can assume $j$ is positive this time as we can assume $k$ and $n$ are positive).
"But this means that n divides m, which again is a contradiction."
Actually it's not a contradiction if $n = 1$.
But this isn't a contradiction that needed to be reached. $k/n$ is a factor of $k$ only makes sense if $n|k$ which would imply $n|m = 2k$.
See this proof that if $n$ is not a perfect square then $\sqrt{n}$ is irrational:
Follow-up Question: Proof of Irrationality of $\sqrt{3}$
The proof starts by saying that if $n$ is not a perfect square then there is a $k$ such that $k^2 < n < (k+1)^2$. The proof breaks down if $k^2 = n$.
Note that this proof does not use divisibility.
• It is misleading to say the proof does not use divisibility. It essentially uses the division algorithm to achieve descent on denominators. I explain this further in this May 20, 2009 sci.math post, where I highlight the beautiful view of irrationality proofs in terms of Dedekind's conductor ideal. – Bill Dubuque Jan 26 '17 at 21:18
• This has nothing to do with the question asked, so it should be a comment, not an answer. – Bill Dubuque Jan 26 '17 at 21:26 | 2019-04-24T13:55:45 | {
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http://15462.courses.cs.cmu.edu/fall2015/lecture/texture/slide_003 | Previous | Next --- Slide 3 of 61
Back to Lecture Thumbnails
BryceSummers
Question: Can anyone use some linear algebra to justify why $R^{-1} = R^{T}$ in this case. (Matrices are not equal to their transpositions in general.) In other words, what is special about matrix $R$?
lucida
R is special because its columns are orthonormal vectors. We know that the dot product of a vector with itself is the norm of that vector. We also know that the dot product of two orthogonal vectors is 0.
We want a matrix $R^{-1}$ such that when multiplied with R we get the identity matrix. Now we observe that if we transpose R, each row i in $R^{T}$ will be column i in $R$. This means when we multiply $R^{T}$ by $R$, the only non zero value in each column i of the resulting matrix will be the value at the jth row where j = i since all other values are the result of a dot product between orthogonal vectors. But when j = i we have the product of two identical vectors of unit length, so their dot product is 1. This gives us the diagonal matrix with 1s along the diagonal, i.e. the identity matrix.
BryceSummers
@lucida, That is a great explanation! | 2020-05-31T02:10:02 | {
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https://math.stackexchange.com/questions/4130618/if-x-y-and-z-are-independent-then-sigmax-y-and-sigmaz-are-indepen/4131328 | # If $X,Y$ and $Z$ are independent, then $\sigma(X,Y)$ and $\sigma(Z)$ are independent
If we have 3 independent random variables $$X,Y$$ and $$Z$$ on a probability space $$(\Omega,\mathscr{F},\mathbb{P})$$, then how do we show that $$\sigma(X,Y)$$ and $$\sigma(Z)$$ are independent?
It feels intuitively obvious, but I'm having real difficulty making it precise. I feel like using the fact that if two measures $$\mu _1$$ and $$\mu _2$$ are equal on a $$\pi$$ system $$\mathscr{A}$$ then they are equal on $$\sigma (\mathscr{A})$$ could be useful, but am unsure exactly what measures/$$\pi$$ systems to use.
Any help would be really appreciated - thanks! :)
• Where did you find this exercise? May 7, 2021 at 21:39
• This came up as a past paper question for an upcoming exam I'm revising for :) May 7, 2021 at 21:44
A similar proof for what follows can be found on page 39 of the book Probability with Martingales (1991) by David Williams. Define the following collections of sets:
\begin{align*}\Pi_1&=\Big\{\lbrace \omega\in\Omega: X(\omega)\le x\rbrace\cap \lbrace \omega\in\Omega: Y(\omega)\le y\rbrace:x,y\in\mathbb{R}\Big\},\\ \Pi_2&=\Big\{\lbrace \omega\in\Omega: Z(w)\le z\rbrace:z\in\mathbb{R}\Big\}. \end{align*}
You can readily show that $$\Pi_1$$ and $$\Pi_2$$ are $$\pi$$-systems such that $$\sigma(\Pi_1)=\sigma(X,Y)\quad\text{and}\quad \sigma(\Pi_2)=\sigma(Z).$$
To prove that $$\sigma(X,Y)$$ and $$\sigma(Z)$$ are independent, by definition we must show that $$P(A\cap B)=P(A)P(B)$$ for all $$A\in \sigma(X,Y)$$ and $$B\in \sigma(Z)$$.
To that end, fix $$A\in \Pi_1$$ and define the following two measures (you can prove that they are measures) $$\mu_{1},\mu_{2}:\mathcal{F}\rightarrow [0,\infty)$$ as follows: $$\mu_1(B)=P(A\cap B)\quad\text{and} \quad \mu_2(B) = P(A)P(B).$$ Since $$A\in \Pi_1$$, we can write $$A=\lbrace X\le x\rbrace\cap \lbrace Y\le y\rbrace$$ for some $$x,y\in\mathbb{R}$$. By independence of the random variables $$X$$, $$Y$$, and $$Z$$, if $$B=\lbrace{Z\le z\rbrace}\in \Pi_2$$ then we find that \begin{align*}\mu_1(B) &= P(A\cap B) \\&= P\left(\lbrace X\le x\rbrace\cap \lbrace Y\le y\rbrace\cap \lbrace{Z\le z\rbrace}\right) \\&= P\left(\lbrace X\le x\rbrace\right)\cdot P\left(\lbrace Y\le y\rbrace\right)\cdot P\left(\lbrace{Z\le z\rbrace}\right)\\&=P\left(\lbrace X\le x\rbrace\cap \lbrace Y\le y\rbrace\right)\cdot P\left(\lbrace{Z\le z\rbrace}\right)\\&=P(A) P(B)\\&=\mu_2(B). \end{align*} This proves that the measures $$\mu_1$$ and $$\mu_2$$ agree on the $$\pi$$-system $$\Pi_2$$. Furthermore, note that $$\mu_1(\Omega)=\mu_2(\Omega) = P(A)<\infty$$. Now we can use the following fact: finite measures that agree on a $$\pi$$-system also agree on the $$\sigma$$-algebra generated by the $$\pi$$-system. Therefore, $$P(A\cap B) = P(A)P(B)\quad\text{for all}\quad A\in \Pi_1\quad\text{and}\quad B\in \sigma(Z).$$ Now fix $$B\in \sigma(Z)$$ and define two more measures $$\mu_{3},\mu_{4}:\mathcal{F}\rightarrow [0,\infty)$$ as follows: $$\mu_3(A)=P(A\cap B)\quad\text{and} \quad \mu_2(A) = P(A)P(B).$$ Our previous argument showed that $$\mu_3$$ and $$\mu_4$$ agree on the $$\pi$$-system $$\Pi_1$$. Furthermore, $$\mu_3(\Omega)=\mu_4(\Omega)=P(B)<\infty$$. Using the same fact as before, we conclude that $$\mu_3$$ and $$\mu_4$$ agree on $$\sigma(X,Y)$$. This completes the proof.
This argument can be generalized as follows: Let $$X_1,X_2,\dots$$ be a sequence of independent random variables. Then $$\sigma(X_1,\dots, X_n)$$ and $$\sigma(X_{n+1}, X_{n+2}, \dots)$$ are independent for each $$n\in\mathbb{N}$$; see page 47 in Williams (1991).
• Could you show me how $\sigma(\Pi_1)=\sigma(X,Y)$? I don't really get how to extend Williams' argument for the individual random variable to two or more. May 8, 2021 at 1:44
• By definition, $\sigma(X,Y)$ is the smallest $\sigma$-algebra that makes both $X$ and $Y$ measurable. Since it (obviously) must contain the sets in $\Pi_1$, we have $\sigma(\Pi_1)\subseteq \sigma(X,Y)$. Conversely, note that $\lbrace X\le a\rbrace = \bigcup_{n=1}^{\infty}(\lbrace X\le a\rbrace\cap\lbrace Y\le n\rbrace\in\sigma(\Pi_1)$ for all $a\in\mathbb{R}$. Similarly, $\lbrace Y\le b\rbrace\in\sigma(\Pi_1)$ for all $b\in\mathbb{R}$. So $\sigma(X)\subseteq \sigma(\Pi_1)$ & $\sigma(Y)\subseteq\sigma(\Pi_1)$. Hence, $\sigma(X,Y)=\sigma\left(\sigma(X)\cup\sigma(Y)\right)\subseteq\sigma(\Pi_1)$. May 8, 2021 at 17:56
• @Snoop If any of these measure-theoretic concepts are unfamiliar, my best advice is invest in an analysis book. Two good options are Real & Complex Analysis by Walter Rudin or Real Analysis: Modern Techniques and Their Applications by Gerald B. Folland. They are written for a graduate student in math and assume basic knowledge from undergraduate analysis (elementary set theory, limits, etc.), but are otherwise self-contained. Even measure-theoretic probability texts usually devote a chapter (or an appendix) to analysis. This type of argument is pretty standard. May 8, 2021 at 18:06 | 2022-08-16T07:32:54 | {
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https://math.stackexchange.com/questions/3213682/let-a-be-a-101-element-subset-of-the-set-s-1-2-ldots-1000000 | Let $A$ be a $101$-element subset of the set $S=\{1,2,\ldots,1000000\}$
Let $$A$$ be a $$101$$-element subset of the set $$S=\{1,2,\ldots,10^6\}$$. For each $$s\in S$$ let $$A_s = A+s = \{a+s \mid a\in A\}$$ Prove that there exist $$B\subset S$$ such that $$|B|=100$$ and the sets in a family $$\{A_b \mid b\in B\}$$ are pairwise disjoint.
Is a following proof corect?
Let $$B$$ be maximal such set that sets in $$\{A_b \mid b\in B\}$$ are pairwise disjoint and let $$|B|=k$$.
Then for each $$b'\in B'= S\setminus B$$ exists such $$b \in B$$ that $$A_b\cap A_{b'}\ne \emptyset$$, i.e. exists $$a_1,a_2\in A$$ so that $$b' = a_1-a_2+b\;\;\;\;\;\;\;\;\;(*)$$
Now make a bipartite graph with $$B'$$ on the left side and on right side: $$C:= \{(a_1,a_2,b)\mid a_1,a_2\in A, a_1\ne a_2, b\in B\}$$
Clearly $$|C| = 101\cdot 100\cdot k$$ and $$|B'| = 10^6-k$$.
Connect $$b'\in B'$$ with $$(a_1,a_2,b)\in C$$ iff $$b' = a_1-a_2+b$$. Then clearly each triple has degree at most $$1$$ and each $$b'$$ has because of $$(*)$$ degree at least $$1$$. By double counting we have: $$10^6-k=|B'|\leq |C| = 101\cdot 100\cdot k$$
so $$k\geq {10^6\over 101\cdot 100+1}>99$$
So $$k\geq 100$$ and we are done.
Apart from a proof verification, I'm interested in probabilistic solution of this problem.
• Can you explain the inequality $|B'|\le|C|$? – W-t-P May 4 at 20:06
• Each vertex in B' has degree at least one and ecah in C at most one – Aqua May 4 at 20:08
1 Answer
The proof is fine.
A more general setup for the proof would be to consider the graph $$G$$ with vertex set $$S$$ and an edge between $$x, y\in S$$ whenever $$A_x \cap A_y = \varnothing$$. (It might be more convenient to put an edge between $$x,y$$ for every element of $$A_x \cap A_y$$, making $$G$$ a multigraph.) Then we are looking for an independent set of $$100$$ vertices in $$G$$.
Your bipartite graph between $$B'$$ and $$C$$ has a sort of incarnation within $$G$$. Consider the subgraph of $$G$$ consisting of all edges between $$B$$ and $$B'$$. Every $$b \in B$$ has degree $$101 \cdot 100$$ in $$G$$ ($$b$$ has an edge to $$b + a_1 - a_2$$ for every $$a_1, a_2 \in A$$ with $$a_1 \ne a_2$$), and these edges must all go to $$B'$$, since $$B$$ is independent. Every $$b' \in B'$$ has at least $$1$$ edge to $$B$$, because $$B$$ is a maximal independent set. So the number of edges between $$B$$ and $$B'$$ is $$101 \cdot 100 \cdot |B|$$, but it's also at least $$|B'| = 10^6-|B|$$. Therefore $$10100 |B| \ge 10^6-|B| \iff |B| \ge \frac{10^6}{10101} > 99.$$ This is essentially a restatement of your argument: rather than have $$10100$$ elements of $$C$$ with degree $$1$$ for every $$b \in B$$, we combine them into a single vertex $$b$$ with degree $$10100$$.
More generally, this shows that in a graph (or multigraph) $$G$$ with $$n$$ vertices and maximum degree $$\Delta(G)$$, there is an independent set of size $$\frac{n}{\Delta(G)+1}$$.
The probabilistic method can be used here to get a bound that's better in general, but not an improvement in this problem. In general, we can get to $$\frac{n}{d+1}$$, where $$d$$ is the average degree in $$G$$. But here, the average degree is also quite close to $$10100$$, so this is not much help.
Here is the probabilistic argument. Randomly permute the elements of $$S$$ as $$b_1, b_2, \dots, b_{10^6}$$, and go through them one at a time to create an independent set $$B$$. For each $$b_i$$, add $$b_i$$ to $$B$$ if every element of the form $$b_i + a_1 - a_2$$ (with $$a_1, a_2 \in A$$ and $$a_1 \ne a_2$$) comes after $$b_i$$ in our random ordering.
This is guaranteed to create an independent set: if $$b_i + a_1 - a_2 = b_j$$, then either $$i (and so we are guaranteed not to have picked $$b_j$$) or $$i>j$$ (and so we are guaranteed not to have picked $$b_i$$). For any $$b \in S$$: of $$b$$ and its at-most-$$10100$$ adjacent elements, each is equally likely to come first in the random ordering, and it's $$b$$ itself with probability $$\frac{1}{10101}$$, in which case we add it to $$B$$.
So the expected size of $$B$$ is at least $$\frac1{10101}|S| = \frac{10^6}{10101} > 99$$, as before, and therefore some random ordering produces a $$B$$ of size at least $$100$$.
• If I'm not mistaken that is basicly method used also in a proof of Caro-Wei Theorem? web.evanchen.cc/handouts/ProbabilisticMethod/… – Aqua May 4 at 21:14
• Yes, that is exactly what's happening. I've always thought of this as Turán's theorem, though. (I guess Turán is the average-degree statement, while Caro-Wei is the stronger form with $\sum_{v \in V} \frac1{\deg(v)+1}$, which gives us a harmonic mean instead.) – Misha Lavrov May 4 at 21:19 | 2019-09-19T08:48:20 | {
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https://math.stackexchange.com/questions/2991201/find-a-general-solution-for-int-0-infty-sin-leftxn-right-dx | # Find a general solution for $\int_{0}^{\infty} \sin\left(x^n\right)\:dx$
So, I was recently working on the Sine Fresnal integral and was curious whether we could generalise for any Real Number, i.e.
$$I = \int_{0}^{\infty} \sin\left(x^n\right)\:dx$$
I have formed a solution that I'm uncomfortable with and was hoping for qualified eyes to have a look over.
So, the approach I took was to employ Complex Numbers (I forget the name(s) of the theorem that allows this).
But
$$\sin\left(x^n\right) = \Im\left[-e^{-ix^n}\right]$$
And so, n
$$I = \int_{0}^{\infty} \sin\left(x^n\right)\:dx = \Im\left[\int_{0}^{\infty} -e^{-ix^n}\:dx \right]= -\Im\left[\int_{0}^{\infty} e^{-\left(i^{\frac{1}{n}}x\right)^{n}}\:dx \right]$$
Applying a change of variable $$u = i^{\frac{1}{n}}x$$ we arrive at:
\begin{align} I &= -\Im\left[i^{-\frac{1}{n}}\int_{0}^{\infty} e^{-u^{n}}\:du \right] \\ &= -\Im\left[i^{-\frac{1}{n}}\frac{\Gamma\left(\frac{1}{n}\right)}{n} \right]\\ &= \sin\left(\frac{\pi}{2n}\right)\frac{\Gamma\left(\frac{1}{n}\right)}{n} \end{align}
My area of concern is in the substitution. As $$i^{-\frac{1}{n}} \in \mathbb{C}$$, I believe the limits of the integral should have been from $$0$$ to $$i^{-\frac{1}{n}}\infty$$. Is that correct or not?
I'm also struggling with bounds on $$n$$ for convergence. Is this expression valid for all $$n\in\mathbb{R}$$
Any guidance would be greatly appreciated
• Where does the negative sign in the imaginary part come from – Darkrai Nov 9 '18 at 10:23
• @Manthanein - I wanted the integral to take the form of the Gamma Function and so to get the negative part in the exponential, I used the property that $e^{-ix} = \cos(x) -sin(x)i$ or $-e^{ix} =-\cos(x) + \sin(x)$ – user150203 Nov 9 '18 at 10:28
• @Manthanein - I just realised I hadn't used that form in the question. Have now corrected. Sorry for any confusion caused. – user150203 Nov 9 '18 at 10:30
• You have indeed moved the contour, so you need to show that the integral from $\infty$ to $i^{-1/n} \infty$ is zero, by an appropriate limiting argument. – Richard Martin Nov 9 '18 at 10:32
• @RichardMartin - Is my result then a consequence of bad maths giving the right answer or is there a link that I'm unaware of? – user150203 Nov 9 '18 at 10:36
Another approach substitutes $$y=x^n$$ and writes $$y^{1/n-1}$$ in terms of a Gamma integral, viz. $$I=\Im\int_0^\infty\frac{1}{n}y^{1/n-1}\exp iy dy=\Im\int_0^\infty\int_0^\infty\frac{1}{n\Gamma(1/n)}z^{-1/n}\exp -y(z-i)dydz.$$By Fubini's theorem, and using $$\Im\frac{1}{z-i}=\frac{1}{1+z^2}$$,$$I=\int_0^\infty\frac{1}{n\Gamma(1/n)}\frac{z^{-1/n}}{1+z^2}dz.$$Then the substitution $$z=\tan u$$ obtains a Beta integral, which can be rewritten in terms of Gamma functions, and the result you've claimed is proven true, by the reflection formula of the Gamma function.
Some Hints:
$$I=\int_0^{\infty} \sin (x^n)dx$$ On substitution $$x^n=t$$ and using the series of $$\sin$$ we get $$I=\frac 1n \int_0^{\infty} t^{\frac 1n} \left(\sum_{k=0}^{\infty} (-1)^k \frac {t^{2k}k!}{(2k+1)!k!} \right) dt$$
On substituting $$t^2=u$$ we get $$I= \frac {1}{2n} \int_0^{\infty} u^{\frac {1-n}{2n}}\left(\sum_{k=0}^{\infty} \frac {\frac {k!}{(2k+1)!}}{k!} (-u)^k \right) du$$
$$I=\frac {1}{2n} \Gamma(s)\phi(-s)$$ where $$\phi(k)=\frac {k!}{(2k+1)!}$$ and $$s=\frac {n+1}{2n}$$
Hence along with properties of Gamma function, Mellin Transform and the Euler's reflection formula we get $$I=\frac {\pi}{2n\cos \left(\frac {\pi}{2n}\right)\Gamma \left(1-\frac 1n\right)}=\sin \left(\frac {\pi}{2n}\right)\frac {\Gamma\left(\frac 1n\right)}{n}$$
With a special case of $$n=2$$ we get the value of special integral popularly known as Fresnel integral with limit as $$x$$ tends to infinity
Start out with a couple of integration by parts: \begin{align} \int_0^\infty\sin(x)\,e^{-xy}\,\mathrm{d}x &=-\frac1y\int_0^\infty\sin(x)\,\mathrm{d}e^{-xy}\tag1\\ &=\frac1y\int_0^\infty\cos(x)\,e^{-xy}\,\mathrm{d}x\tag2\\ &=-\frac1{y^2}\int_0^\infty\cos(x)\,\mathrm{d}e^{-xy}\tag3\\ &=\frac1{y^2}-\frac1{y^2}\int_0^\infty\sin(x)\,e^{-xy}\,\mathrm{d}x\tag4\\ &=\frac1{y^2+1}\tag5 \end{align} Explanation:
$$(1)$$: prepare to integrate by parts
$$(2)$$: integrate by parts
$$(3)$$: prepare to integrate by parts
$$(4)$$: integrate by parts
$$(5)$$: add $$\frac{y^2}{y^2+1}$$ times $$(4)$$ to $$\frac1{y^2+1}$$ times the LHS of $$(1)$$
Now write \begin{align} \int_0^\infty\sin\left(x^n\right)\,\mathrm{d}x &=\frac1n\int_0^\infty\sin(x)\,x^{\frac1n-1}\,\mathrm{d}x\tag6\\[3pt] &=\frac1{n\,\Gamma\!\left(1-\frac1n\right)}\int_0^\infty\sin(x)\int_0^\infty y^{-\frac1n}e^{-xy}\,\mathrm{d}y\,\mathrm{d}x\tag7\\ &=\frac1{n\,\Gamma\!\left(1-\frac1n\right)}\int_0^\infty y^{-\frac1n}\int_0^\infty\sin(x)\,e^{-xy}\,\mathrm{d}x\,\mathrm{d}y\tag8\\ &=\frac1{n\,\color{#C00}{\Gamma\!\left(1-\frac1n\right)}}\color{#090}{\int_0^\infty\frac{y^{-\frac1n}}{y^2+1}\,\mathrm{d}y}\tag9\\ &=\color{#C00}{\frac{\Gamma\!\left(\frac1n\right)\sin(\frac\pi{n})}{\color{#000}{n}\pi}}\color{#090}{\frac\pi2\sec\left(\frac\pi{2n}\right)}\tag{10}\\[9pt] &=\Gamma\!\left(1+\frac1n\right)\sin\left(\frac\pi{2n}\right)\tag{11} \end{align} Explanation:
$$\phantom{1}(6)$$: substitute $$x\mapsto x^{1/n}$$
$$\phantom{1}(7)$$: $$\int_0^\infty y^{-\frac1n}e^{-xy}\,\mathrm{d}y=x^{\frac1n-1}\Gamma\!\left(1-\frac1n\right)$$
$$\phantom{1}(8)$$: Fubini
$$\phantom{1}(9)$$: apply $$(5)$$
$$(10)$$: $$(4)$$ from this answer for the green, and $$(2)$$ from the same answer for the red
$$(11)$$: simplify
Here is an alternative approach that avoids complex numbers and series altogether. To get round these two obstacles I will use a property of the Laplace transform.
Let $$I = \int_0^\infty \sin (x^n) \, dx, \qquad n > 1.$$ We begin by enforcing a substitution of $$x \mapsto x^{1/n}$$. This gives $$I = \frac{1}{n} \int_0^\infty \frac{\sin x}{x^{1 - 1/n}} \, dx.$$
The following useful property (does this result have a name? It would be so much nicer if it did!) for the Laplace transform will be used: $$\int_0^\infty f(x) g(x) \, dx = \int_0^\infty \mathcal{L} \{f(x)\} (t) \cdot \mathcal{L}^{-1} \{g(x)\} (t) \, dt.$$ Noting that $$\mathcal{L} \{\sin x\}(t) = \frac{1}{1 + t^2},$$ and $$\mathcal{L}^{-1} \left \{\frac{1}{x^{1-1/n}} \right \} (t)= \frac{1}{\Gamma (1 - \frac{1}{n})} \mathcal{L}^{-1} \left \{\frac{\Gamma (1 - \frac{1}{n})}{x^{1-1/n}} \right \} (t) = \frac{t^{-1/n}}{\Gamma (1 - \frac{1}{n})},$$ then \begin{align} I &= \frac{1}{n} \int_0^\infty \sin x \cdot \frac{1}{x^{1 - \frac{1}{n}}} \, dx\\ &= \frac{1}{n} \int_0^\infty \mathcal{L} \{\sin x\} (t) \cdot \mathcal{L}^{-1} \left \{\frac{1}{x^{1 - \frac{1}{n}}} \right \} (t) \, dt\\ &= \frac{1}{n\Gamma (1 - \frac{1}{n})} \int_0^\infty \frac{t^{-1/n}}{1 + t^2} \, dt. \end{align} Enforcing a substitution of $$t \mapsto \sqrt{t}$$ yields \begin{align} I &= \frac{1}{2 n \Gamma \left (1 - \frac{1}{n} \right )} \int_0^\infty \frac{t^{-\frac{1}{2} - \frac{1}{2n}}}{t + 1} \, dt\\ &= \frac{1}{2 n \Gamma \left (1 - \frac{1}{n} \right )} \operatorname{B} \left (\frac{1}{2} - \frac{1}{2n}, \frac{1}{2} + \frac{1}{2n} \right )\\ &= \frac{1}{2 n \Gamma \left (1 - \frac{1}{n} \right )} \Gamma \left (\frac{1}{2} - \frac{1}{2n} \right ) \Gamma \left (\frac{1}{2} + \frac{1}{2n} \right ). \tag1 \end{align} Applying Euler's reflexion formula we have $$\Gamma \left (\frac{1}{2} - \frac{1}{2n} \right ) \Gamma \left (\frac{1}{2} + \frac{1}{2n} \right ) = \frac{\pi}{\sin \left (\frac{\pi}{2n} + \frac{\pi}{2} \right )} = \frac{\pi}{\cos \left (\frac{\pi}{2n} \right )},$$ and $$\Gamma \left (1 - \frac{1}{n} \right ) = \frac{\pi}{\sin \left (\frac{\pi}{n} \right ) \Gamma \left (\frac{1}{n} \right )}.$$ So (1) becomes $$I = \frac{\sin (\frac{\pi}{n} ) \Gamma (\frac{1}{n})}{2n \cos (\frac{\pi}{2n} )},$$ or $$I = \sin \left (\frac{\pi}{2n} \right ) \frac{\Gamma \left (\frac{1}{n} \right )}{n}, \qquad n > 1$$ where in the last line the double angle formula for sine has been used. | 2019-04-21T03:14:20 | {
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https://math.stackexchange.com/questions/3164773/does-the-definition-of-compactness-require-the-open-cover-to-consist-of-subsets | # Does the definition of compactness require the open cover to consist of subsets?
Sorry if the title is a bit unclear, but I'm stuck on the definition of compactness for metric spaces using open covers. So our professor wrote (word for word) that in a metric space $$(X,d)$$, an open cover is a family $$\{U_i\}_{i\in I}$$ with
$$U_i\subset X$$ open,
$$\bigcup_{i\in I}U_i=X$$
and that $$X$$ is compact if from every open cover $$\{U_i\}_{i\in I}$$ finitely many $$U_{i_1},...,U_{i_r}$$ can be chosen such that $$U_{i_1}\cup U_{i_2} \cup ... \cup U_{i_r}=X$$.
Is this definition correct? Because in one of our exercises, they ask us to find an open cover (that doesn't have a finite subcover) for $$X=[0,1]\cap \Bbb Q$$ with the Euclidian distance, but I can't for the life of me think of any set that is a subset of $$X$$, contains $$0$$ or $$1$$, and on top of that is an open set - those three things seem mutually exclusive to me. Was our prof mistaken in writing that $$U_i$$ is a subset of $$X$$ in the definition for open covers?
Any help would be greatly appreciated!
• Your prof was not mistaken. It might help to recall the definition of the subspace topology
– jgon
Mar 27 '19 at 16:56
Open means relatively open. A subset of $$U$$ of $$X = [0,1] \cap \mathbf Q$$ is open if there is an open set $$O \subset \mathbf R$$ satisfying $$U = O \cap X$$.
A subset $$K$$ of a metric space $$X$$ is compact if and only if for every any collection of open sets $$U_i\subseteq X,i\in I$$ for which $$K\subseteq \bigcup_{i\in I} U_i$$, there is a finite set $$F\subseteq I$$ for which $$K\subseteq \bigcup_{i\in F}U_i$$. Note that $$U_i$$ need only be subsets of $$X$$.
However, $$K$$ being compact as a subset of $$X$$ is equivalent to $$(K,d|_K)$$ being a compact space. Here, $$d|_K$$ is the metric on $$X$$ with its domain restricted to $$K\times K$$. It can be shown that $$V\subseteq K$$ is open in $$(K,d|_K)$$ if and only if $$V=U\cap K$$ for some open $$U\subseteq X$$.
For your problem, for example, letting $$K=[0,1]\cap \mathbb Q$$, then $$(a,b)\cap [0,1]\cap \mathbb Q$$ would be an open set for any $$a,b\in \mathbb R$$. Even though such sets are not open in $$\mathbb R$$, they are open when viewed as subsets of $$K$$ with the metric inherited from $$\mathbb R$$.
e.g. $$[0,\frac12)\cap \mathbb{Q}$$ is an open set of $$X=[0,1]\cap \mathbb{Q}$$ in its subspace topology. | 2022-01-16T10:10:49 | {
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http://math.stackexchange.com/questions/387657/choosing-elements-from-sets | # Choosing elements from sets
OK, so I've always been terrible at combinatorics and I'm trying to generalize some combinatorial problems and I can't figure out where I'm going wrong. Take the following problem:
Assume we are given disjoint finite sets $A_1,A_2,\ldots,A_n$ with $A=\bigcup_i A_i$. How many sets can we form with exactly one element chosen from each $A_i$?
It would seem that the answer is $\prod_i |A_i|$. But then take the following problem:
Assume we draw 4 cards out of a deck. What is the probability that they are of different suits?
This would have the obvious solution:
$$\frac{52\times 39\times 26\times 13}{{52\choose 4}}$$
However, if I let $A_1,\ldots,A_4$ be the different suits, then the above would give me:
$$\frac{13^4}{{52\choose 4}}.$$
Where's my logic going wrong?
I had originally reasoned that if we write $\mathcal{F}=\{(x_1,\ldots,x_n)\mid x_i\in A_i \}$, then we have a map:
$$f:\mathcal{F}\to \mathcal{P}(A),\;\;f(x_1,\ldots,x_n)=\{x_1,\ldots,x_n\}$$
This is clearly a bijection to the subset of elements of $\mathcal{P}(A)$ that satisfy the requirement. Thus, it should be the right answer. My book, however, gave the following as the correct answer to the card problem:
$$\frac{52\times 39\times 26\times 13}{{52\choose 4}}$$
while based on the comments what I thought was the right answer:
$$\frac{13^4}{{52\choose 4}} = \frac{52\times 39\times 26\times 13}{52\times 51\times 50\times 49}$$
is actually correct.
-
I guess the $A_i$ are disjoint and you want to choose exactly one element from each $A_i$? – Michael Greinecker May 10 '13 at 15:10
Yes, the $A_i$ are disjoint. – user75789 May 10 '13 at 15:12
It seems like based on the comments, the second answer is correct (this is what I initially wrote down). However, the solution sheet in the back of my textbook claims that the first one is correct... I was puzzled by this, since there's an obvious bijection from the set of size $13^4$ to the subsets of size $4$ of the deck that satisfies the property. – user75789 May 10 '13 at 15:32
Suppose that the deck has been separated into piles, each containing the cards of one suit, and you’re counting the number of ways to draw one card from each pile. There are $13^4$ ways to do that, since you may take any of the $13$ cards in each of the $4$ piles. Every one of these ways gives you a card from each suit, so in this scheme the probability of drawing one card from each suit is
$$\frac{13^4}{13^4}=1\;.$$
Note that the denominator is not $\binom{52}4$: in this scheme we are not selecting $4$ cards from a single pile of $52$ cards, but rather one card from each of $4$ piles of $13$ cards.
Now return to the original scheme: the suits have not been separated, and you’re drawing $4$ cards from a single pile of $52$ cards. We’ll count the ways to get one card from each suit. The first card drawn may be from any suit, so there are $52$ possible choices. The second must be from one of the other three suits, so there are $52-13=39$ possible choices. Similarly, there are $26$ choices for the third card and $13$ for the fourth card, for a total of $52\cdot39\cdot26\cdot13$ possible sequences of draws resulting in one card from each suit. This calculation uses exactly the same multiplication principle as the calculation that there are $13^4$ ways to draw one card from each suit in the first scheme; it’s just that the relevant sets $A_i$ are different.
Your denominator of $\binom{52}4$, however, is incorrect: it’s the number of sets of $4$ cards that can be drawn from the $52$-card deck, and in the numerator you’re counting $4$-card sequences, not sets. There are actually $52\cdot51\cdot50\cdot49$ possible sequences of draws: the first card can be any of the $52$, the second can be any of the remaining $51$, and so on. Thus, in this scheme the probability of getting one card from each suit is
$$\frac{52\cdot39\cdot26\cdot13}{52\cdot51\cdot50\cdot49}\;.$$
The denominator here is actually $4!\binom{52}4$, the extra factor of $4!$ accounting for the fact that each set of $4$ cards can be drawn in any of $4!$ different orders and therefore corresponds to $24!$ different $4$-card sequences.
-
What bugs me about combinatorics is that in our class we had 10 students and usually about 5 different answers to each question. How do I make this stuff rigorous, since at each step there's a lot of appeal to intuition? I'm curious, since the first answer to the card problem gives the same numeric answer as the solution in the back of the textbook... – user75789 May 10 '13 at 15:27
@user75789: It is rigorous: there’s no intuition involved. Experience in analyzing such problems, perhaps, but not intuition. I’m not sure which answer you’re calling the first one; the only correct answer appearing either in your question or my answer is $$\frac{52\cdot39\cdot26\cdot13}{52\cdot51\cdot50\cdot49}\;.$$ – Brian M. Scott May 10 '13 at 15:30
The second answer is the correct one. which is also the same as the direct calculation: $$\frac{52}{52}\frac{39}{51}\frac{26}{50}\frac{13}{49}$$
The first calculation is obviously wrong, as the result is greater than $1$. The logical explanation, is that the way you count the possible number of "good" draws is ordered, e.g. the same set of $4$ different suited cards is counted $4!$ times. However, is the denominator, you counted only unordered draws. Divide the first answer by $4!$ and you will get the same answer is your second, which is right.
- | 2015-01-28T12:04:37 | {
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https://math.stackexchange.com/questions/1400217/proving-that-a-function-is-odd/1400255 | # Proving that a function is odd
Assume that there exists a function $f:\mathbb{R}\to\mathbb{R}$ that is bijective and satisfies $$f(x) + f^{-1}(x)=x$$ for all $x$. Here $f^{-1}$ is the inverse function. Show that $f$ is odd.
This was a brain-teaser given to me by a friend. Two other related questions are:
1. Show that $f$ is discontinuous
2. Give an example of such a function (if indeed one exists).
Edit: As an initial idea, maybe approaching the problem graphically would help? A function and its inverse are reflections of each other about $y=x$ on the $x$-$y$ plane. Does this lead to anywhere?
• I assume the domain must be $\mathbb R$? – 5xum Aug 17 '15 at 12:53
• Yes the function is from the Reals to the Reals.\ – Iconoclast Aug 17 '15 at 13:08
• One can show that $f(0)=0$ as follows: Let $f(0)=a$, then $f(0)+f^{-1}(0)=0$ or that $f^{-1}(0)=-a$. Therefore, $f(-a)=0$ and $f(0)=a$. Now, $f(-a)+f^{-1}(-a)=-a$, but since $f(-a)=0$, $f^{-1}(-a)=-a$ or that $f(-a)=-a$ so $-a=0$. – Michael Burr Aug 17 '15 at 13:10
Here is finally a constructive example of a solution, continuous except on a countable set.
Let $\phi = \frac{1+\sqrt5}{2}$, and let $$f(x) = \begin{cases} 0 & \text{if }x = 0 \\ -\phi x & \text{if }|x| \in [\phi^{3k}, \phi^{3k+1}), k \in \mathbb Z \\ \phi x & \text{if } |x| \in [\phi^{3k+1}, \phi^{3k+2}), k \in \mathbb Z \\ \phi^{-2} x & \text{if } |x| \in [\phi^{3k+2}, \phi^{3k+3}), k \in \mathbb Z \\ \end{cases}$$
Then $f$ is a bijection, whose inverse is $$f^{-1}(x) = \begin{cases} 0 & \text{if }x = 0 \\ \phi^2 x & \text{if }|x| \in [\phi^{3k}, \phi^{3k+1}), k \in \mathbb Z \\ -\phi^{-1} x & \text{if } |x| \in [\phi^{3k+1}, \phi^{3k+2}), k \in \mathbb Z \\ \phi^{-1} x & \text{if } |x| \in [\phi^{3k+2}, \phi^{3k+3}), k \in \mathbb Z \\ \end{cases}$$ Checking each case shows that $f(x)+f^{-1}(x) = x$, as required.
• What intuition did you follow to arrive at this result? – templatetypedef Aug 18 '15 at 2:26
• @templatetypedef I tried to build the function on $[1, \lambda)$ by setting $f(x)=-\alpha x$. See my other answer to see why this leads to a cycle $\pm[1, \lambda) \to \pm[\alpha, \alpha\lambda) \to \pm[\alpha+1, (\alpha+1)\lambda)$. Letting $\alpha=\lambda$ and $\alpha^2=\alpha+1$ enables the three intervals to be contiguous. At this point we have built the bijection on $[1,\phi^3)$. We can restart from $\phi^3$ to build the bijection on $[\phi^3, \phi^6)$, and so on. – yoann Aug 18 '15 at 10:40
• A "golden" idea.... – Piquito Aug 30 '15 at 11:56
By plugging in $x=f(y)$ we obtain: $$f(f(y))=f(y)-y$$ Call this assertion $P(y)$ and let $f^{(n)}(x)=\underbrace{f(f(…f(x)...))}_{n \space times}$. Now we have: $$P(f(x)): f^{(3)}(x)=f(f(x))-f(x)=f(x)-x-f(x)=-x \iff f^{(4)}(x)=f(-x)$$ But: $$P\left(f(f(x))\right): f^{(4)}(x)=f^{(3)}(x)-f^{(2)}(x)=-x-(f(x)-x)=-f(x)$$ By combining these equations, we obtain $f(-x)=-f(x)$ and therefore, $f$ is odd.
• Nicely done. but what about continuity and existence? – Iconoclast Aug 17 '15 at 13:30
• Thank you. I'm working on it :) – Redundant Aunt Aug 17 '15 at 13:33
Proof that the function is odd (if it exists):
Suppose that $f(b)=a$. Now, consider $$f(b)+f^{-1}(b)=b.$$ Then, $f^{-1}(b)=b-a$ or that $f(b-a)=b$. Now, consider $$f(b-a)+f^{-1}(b-a)=b-a.$$ By substitution, we have $b+f^{-1}(b-a)=b-a$ or that $f^{-1}(b-a)=-a$. Therefore, $f(-a)=b-a$. Now, consider $$f(-a)+f^{-1}(-a)=-a.$$ By substitution, $b-a+f^{-1}(-a)=-a$ or that $f^{-1}(-a)=-b$. In other words, $f(-b)=-a$, so the function is odd.
• Better than accepted answer. Did you use bijective anywhere? – BCLC Aug 18 '15 at 12:56
• Btw what exactly are b and a? I suspect that that is where bijection is used – BCLC Aug 18 '15 at 12:58
• I only used the bijection to get that $f^{-1}$ exists. In this answer, $a$ and $b$ are arbitrary (satisfying the condition that $f(b)=a$). (The accepted answer is quite clean and provides insight on how to solve other similar questions.) – Michael Burr Aug 18 '15 at 13:03
• googles more on bijection Strange. It seems that the fact of declaring the existence of $f^{-1}$ should be sufficient to mean that the function is bijective. Anyway, you mean repeated compositioning can solve other similar questions? – BCLC Aug 18 '15 at 13:08
First, as others have noted, we must have $f(0) = 0$. Let us now "construct" such a function on $\mathbb R^{*}$ using the Axiom of Choice.
Using $f(x) = x - f^{-1}(x)$, it is easy to see that, if $A \in \mathbb R^{*}$ is such that $f(A) = B$, then the following cycle must occur: $$A \mapsto B \mapsto B - A \mapsto -A \mapsto -B\mapsto A-B \mapsto A$$ (Note that this proves that the function is odd.)
Therefore, an idea to construct such a function is for instance to set $B = -2A$, so that we have the following cycle: $$A \mapsto -2A \mapsto -3A \mapsto -A \mapsto 2A \mapsto 3A \mapsto A$$
If we could partition $\mathbb R^{*}$ into sets of the form $\{-3A,-2A,-A,A,2A,3A\}$, we could easily define the function on each of these sets using the cycle above.
To do so, we can use the Axiom of Choice and take the quotient of $\mathbb R^{*}$ by the equivalence relation $$A \equiv B \Leftrightarrow \exists (p,q) \in \mathbb Z^2, A = \pm 2^p 3^q B$$
The Axiom of Choice gives us a set $X$ such that $\mathbb R^{*} = \coprod_{A \in X} \bigcup_{(p,q)\in\mathbb Z^2}\{\pm2^p3^qA\}$.
Then, for each $A\in X$, we can construct $f$ on $E_A = \bigcup_{(p,q)\in\mathbb Z^2}\{\pm2^p3^qA\}$. Indeed, since $\mathbb Z^2$ is countable, we can order it and use induction to build our cycles. When considering the $n$-th element of $\mathbb Z^2$, either it does not appear in a cycle formed by one of the previous elements of $\mathbb Z^2$, in which case we build $f$ on the cycle formed by the current element; or it does, in which case we skip the element.
It is easy to check that the function built this way has the required properties.
Note that I have no idea if we can build such a function without using the Axiom of Choice, but I doubt so.
Edit: I finally found a constructive solution, see my other answer. (If it is inappropriate to double-post, I will merge both answers.)
• could you kindly tell me what is $\mathbb R^{*}$? Sorry for the rather basic question. – Iconoclast Aug 17 '15 at 17:10
• @Iconoclast Just a notation: $\mathbb R^{*} = \mathbb R \setminus \{0\}$ – yoann Aug 17 '15 at 17:12
• In my opinion it's nice to have your two answers separate. – user21820 Aug 18 '15 at 17:16
If $f$ were continuous, then it would have to be strictly monotonic, since it's bijective. Suppose $f$ is decreasing, $x>0$. Then, since $f(0)=0$ we get $f(x)<0$ and $f^{-1}(x)<0$, but $x=f(x)+f^{-1}(x)$. So $f$ is increasing and $f(x)>0$ when $x>0$. But then, $f^{-1}(x)$ is also positive for $x>0$, so $f(x)<x$. Since $f$ is increasing, $f^{-1}(x)>x$. But that's not possible.
• So does this mean that $f(x)$ is discontinuous everywhere? i.e., continuous on no internal? – Iconoclast Aug 17 '15 at 13:49
• This only shows that it has a discontinuity. I suspect such a function ( if one exists ), is very discontinuous, but I'm still thinking about it. – Callus Aug 17 '15 at 13:56
Let's call $(E)$ the equation $f(y) + f^{-1} (y) = y$
Let $a \in \mathbb{R}$. Since $f$ is bijective, there exists $x \in \mathbb{R}$ such that $f(x) = a$. Now, since $f^{-1}(a) = x$, we have (using $(E)$ with $y = a$) $f(a) = a - x$. Now, using $(E)$ with $y = x$:
$$f^{-1}(x) = x - a$$ Applying $f$ on both side: $$f(x-a) = x$$ Using $(E)$ with $y = x-a$: $$f^{-1}(x-a) = x - a - f(x -a) = x -a -x = a$$ Applying $f$ on both side: $$f (-a ) = x -a$$.
Thus $f(a) = a - x = -f(-a)$. This is true for all $a \in \mathbb{R}$: $f$ is therefore odd. | 2019-07-17T21:16:16 | {
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http://mathhelpforum.com/pre-calculus/170909-show-p-q-th-term-series.html | # Thread: Show that p/q is the ...th term of the series...
1. ## Show that p/q is the ...th term of the series...
Positive rational numbers may be arranged in the form of a simple series: $\frac{1}{1}, \frac{2}{1}, \frac{1}{2}, \frac{3}{1}, \frac{2}{2}, \frac{1}{3}, \frac{4}{1}, \frac{3}{2}, \frac{2}{3}, \frac{1}{4}, ...$ Show that $\frac{p}{q}$ is the $[\frac{1}{2}(p+q-1)(p+q-2)+q]^t^h$ term of the series.
I have already considered arithmetic sequences for p and q but I don't think it works. How should I approach the problem?
Thanks.
2. This expression of p and q is called Cantor pairing function. It maps pair of positive integers onto positive integers in a one-to-one manner. It is useful because it demonstrates that there are as many positive integers as pairs of positive integers.
To understand why it works, consider this picture. The value of the pairing function is shown near each point. Arrow show the direction in which the value of the function increases.
Note that p + q is constant on the diagonals. Let's look at the blue diagonal, where p + q = 5. How many points come before, i.e., how many points are in the green triangle? It's 1 + 2 + 3, or, in terms of p and q, it's 1 + 2 + ... + (p + q) - 2 = 1/2 [(p + q - 2) * (p + q - 1)] (as the sum of an arithmetic progression). If we add the vertical coordinate q, we get the value of the function on the blue diagonal.
3. Thank you for your solution. I didn't expect an advanced topic like the Cantor pairing function to be required for this problem. I was told by a teaching assistant that this problem can be solved by mathematical induction. But he said it may not work because after p or q reaches the end of its arithmetic progression, it "jumps" to a different number. In the inductive step, we assume the statement is true for some n, then we try to show that (p-1)/(q+1) is the (n+1)th term of the series. This would result in the wrong numbers in some cases. Can induction still be used anyway? In a different way maybe?
4. @ akyng
could you post the exact problem. or a copy of it. there is some differences from what you posted and emakarovs expression.
akyng
$[\frac{1}{2}(p+q-1)(p+q-2)+q]^t^h$
emakarov
$1 + 2 + ... + (p + q) - 2 = \frac{1}{2} [(p + q - 2)(p + q - 1)]$
notice you have 1/2 fraction in different grouping, there is +q at the end of your expression and emakarov has a (p+q)-2 nth term. that would give you 0,1,2,3,4 etc, for the sequence.
5. I think our expressions are exactly the same. The green triangle has $1 + 2 + ... + ((p + q) - 2) = \frac{1}{2} [(p + q - 2)(p + q - 1)]$ points. To this one adds q: the position of the current point on the blue diagonal.
6. ah I see
so what would you put for first element? ((2+1)-2)?
7. The first element of the series is 1/1 or 0.5*(1+1-2)(1+1-1)+1. Does someone know something to the question I asked in post #3, about whether mathematical induction can be used for this problem?
8. Does someone know something to the question I asked in post #3, about whether mathematical induction can be used for this problem?
OK, first a couple of definitions.
Let $f(p,q)=(p+q-1)(p+q-2)/2+q$. Also, let $p_n,q_n$ be defined by mutual recursion as follows: $p_1=q_1=1$,
$p_{n+1}=
\begin{cases}
p_n-1 & p_n>1\\
q_n+1 & \mbox{otherwise}
\end{cases}$
and $q_{n+1}=\begin{cases}
q_n+1 & p_n>1\\
1 & \mbox{otherwise}
\end{cases}$
.
The original problem:
Show that $\frac{p}{q}$ is the $[\frac{1}{2}(p+q-1)(p+q-2)+q]^t^h$ term of the series.
asks to prove that for all $p,q$ it is the case that $p=p_{f(p,q)}$, $q=q_{f(p,q)}$. This statement can be broken into two parts.
Proposition 1. For all $p,q$ there exists an $n$ such that $p=p_n$ and $q=q_n$.
Proposition 2. $f(p_n,q_n)=n$ for all $n\ge 1$.
This is proved by induction on $n$. In the induction step, one has to consider two cases depending on whether $p_n>1$.
Then, given $p$ and $q$, we use Proposition 1 to find an $n$ such that $p=p_n$ and $q=q_n$. By Proposition 2, $n=f(p,q)$, which proves the required claim.
9. this one is interesting. I have been doing some of the induction proofs in pre-calculus book but not with 2 variables. so I take a crack at it for my own personal fun.
we start with the chart emakarov posted. we pair the outside numbers as $\frac{p}{q}$. now set $(p+q)=n$. call these elements $a_n$ in a statement involving n. we can fallow the little place numbers to build our $a_n$ values.
now if you run $a_n$ through $((p+q)-2)$ you get a repeat value for each diagonal row like in the chart. if you run $a_n$ through $\frac{1}{2} [(p + q - 2)(p + q - 1)]$ you get the sum of the row values. so I assume that is what we want. so therefore $S_n$ is the nth partial sum of the row values.
$a_n= ((p+q)-2):a_1=0,a_{2-3}=1,a_{4-6}=2,a_{7-10}=3,a_{11-15}=4...$
$S_n= \frac{1}{2} [(p + q - 2)(p + q - 1)]:S_1=0,S_{2-3}=1,S_{4-6}=3,S_{7-10}=6,S_{11-15}=10...$
$a_1=0$ is equal to $S_1=0$ so we have proof that $a_1$ is true. the truth of $a_n$ implies the truth of $a_{n+1}$. we assume that $a_n$ is true and build the $a_{n+1}$ statement.
\begin{aligned}
a_{n+1} &=(p+q)-2) \\
&= (((p+q)+1)-2) \\
&= ((p+q)-1).
\end{aligned}
now we add $a_{n+1}$ to $S_n$.
\begin{aligned}
S_{n+1}
&= \frac{1}{2} [(p + q - 2)(p + q - 1)]+(p+q-1) \\
&= \frac{1}{2} [(p + q - 2)(p + q - 1)+2(p+q-1)] \\
&= \frac{1}{2} [(p + q - 2)(p + q - 1) + 2p + 2q - 2] \\
&= \frac{1}{2} [p^2 + 2pq + q^2 - p - q] \\
&= \frac{1}{2} [(p + q)^2 - p - q] \\
&= \frac{1}{2} [(p + q)^2 - (p + q)] \\
\end{aligned} | 2018-02-22T19:07:30 | {
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http://math.stackexchange.com/questions/145242/finding-matrix-representation/145245 | # Finding Matrix Representation
Problem: Let T: $\mathbb{R}^3\rightarrow\mathbb{R}^3$ be a linear map given by
$$T\left[ \begin{matrix} x\\y\\ z\end{matrix} \right]= \left[ \begin{matrix} 3x-y\\z-x\\z-y\\\end{matrix} \right]$$
1. Find the Matrix representation of T with respect to the canonical basis of $\mathbb{R}^3$, and call it A.
I am not sure how this works. So the cananical basis of $\mathbb{R}^3$ is ${(1,0,0),(0,1,0),(0,0,1)}$ But I am unsure how to get a matrix represenation from a linear operator. Any help is appreciated.
-
Remember how a matrix is defined. The $j$th column tells you where the $j$th basis vector goes. So the first column should be the image of $(1,0,0)$ under $T$, which is $(3,-1,0)$. You can continue like this – Daniel Freedman May 14 '12 at 22:26
Ok great! Thank you guys so much it is so much clearer now. – Mathstudent May 14 '12 at 22:37
I will write the linear map $T$ as $$T(x,y,z)=(3x-y,z-x,z-y).$$ Then we have $$T(1,0,0)=(3,-1,0),$$ $$T(0,1,0)=(-1,0,-1),$$ $$T(0,0,1)=(0,1,1).$$ Then the matrix representation of $T$ is given by $$\left[ \begin{array}{ccc} 3 & -1 & 0 \\ -1 & 0 & 1 \\ 0 & -1 & 1 \\ \end{array} \right].$$ The way to get it is: the vector $T(1,0,0)=(3,-1,0)$ becomes the first column, the vector $T(0,1,0)=(-1,0,-1)$ becomes the second column, and so forth. My linear algebra teacher always says, "Put it to be vertical, put it to be vertical,...".
The advantage of the matrix representation is that; for example if I want to find $T(1,-2,0)$, then I can do it by $$\left[ \begin{array}{ccc} 3 & -1 & 0 \\ -1 & 0 & 1 \\ 0 & -1 & 1 \\ \end{array} \right]\left[ \begin{array}{c} 1 \\ -2 \\ 0 \\ \end{array} \right]=\left[ \begin{array}{c} 5 \\ -1 \\ 2 \\ \end{array} \right],$$ that is, $T(1,-2,0)=(5,-1,2)$.
-
wait shouldn't the matrix represenation be $$\left[ \begin{array}{ccc} 3 & -1 & 0 \\ -1 & 0 & -1 \\ 0 & 1 & 1 \\ \end{array} \right].$$ – Mathstudent May 14 '12 at 22:44
Let $e_1=(1,0,0)$, $e_2=(0,1,0)$, and $e_3=(0,0,1)$.
From the definition of $T$: $$T(e_1) =\left[\matrix{3\cr -1\cr 0 }\right],\quad T(e_2) =\left[\matrix{-1\cr 0\cr -1 }\right],\quad T(e_3) =\left[\matrix{0\cr 1\cr 1 }\right].$$ For an arbitrary vector $v=\left[\matrix{a\cr b\cr c }\right] =a e_1+b e_2+c e_3$, we have using the linearity of $T$ \eqalign{ T(v) &=T(a e_1+b e_2+c e_3)\cr &=a T(e_1)+b T(e_2)+c T(e_3)\cr &=a \left[\matrix{3\cr -1\cr 0 }\right]+ b\left[\matrix{-1\cr 0\cr -1 }\right]+ c\left[\matrix{0\cr 1\cr 1 }\right]\cr &= \left[\matrix{3a-1\cdot b+0\cdot c\cr -1\cdot a+0\cdot b+1\cdot c\cr 0\cdot a+1\cdot -b+1\cdot c }\right]\cr &= \left[\matrix{3&-1&0 \cr -1&0& 1\cr 0&-1& 1}\right] \left[\matrix{a\cr b\cr c}\right].\cr } So the matrix representation of $T$ is $\left[\matrix{3&-1&0 \cr -1&0&1\cr 0&-1&1}\right]$.
Note that the $i^{\rm th}$ column of the matrix representing $T$ is the vector $T(e_i)$. So, we can get the desired!
- | 2015-06-29T23:36:04 | {
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https://dsp.stackexchange.com/questions/18201/recursive-version-of-dft-as-presented-in-cooley-tukey-paper/18798 | # Recursive version of DFT as presented in Cooley-Tukey paper
The seminal paper of Cooley and Tukey provides an iterative method for computing the DFT for a sequence of length $N$. Specifically, they mention a method which utilizes the fact that $N$ can be written as $N = r_1 r_2$. The two intermediate results during computations themselves are DFT computations, therefore, their computation in turn can be recursively estimated.
Is there a reference/method of recursive formulation for computing DFT using the $N=r_1 r_2$ factorization ? The versions of FFT I come across use factorizations where $r_1 = \frac{N}{2}$ and $r_2=2$. I would like to know if there are recursive formulations for general factorization of $N$.
• @MattL. The versions of FFT I come across use factorizations where $r_1 = \frac{N}{2}$ and $r_2=2$. I would like to know if there are recursive formulations for general factorization of $N$. Sep 13, 2014 at 8:54
• The Cooley-Tukey algorithm is not restricted to dividing the transforms into 2 DFTs of length $N/2$. This is just the most popular form. It works for any composite size $N=r_1r_2$. A lot of methods with different factorizations of $N$ are implemented in the FFTW library. You can find more details here. Sep 13, 2014 at 9:11
• it might be interesting if someone here spelled out how to do it in the case where $r_1 \ne \frac{N}{2}$. Sep 13, 2014 at 16:11
• @robertbristow-johnson I finally found the time to do so. Oct 23, 2014 at 11:48
• @MattL., an up vote for that. Oct 26, 2014 at 11:35
The Cooley Tukey method allows general factorizations $N=N_1N_2$. The easiest way to see this is to use index mapping. Consider the length-$N$ DFT of a sequence $x[n]$:
$$X[k]=\sum_{n=0}^{N-1}x[n]W_N^{nk}\tag{1}$$
with $W_N=e^{-j2\pi/N}$. Let's assume that $N$ can be factored as $N=N_1N_2$. Now use the following index mapping
$$n=n_1N_2+n_2,\quad n_1=0,1,\ldots, N_1-1,\quad n_2=0,1,\ldots,N_2-1\\ k=k_1+N_1k_2,\quad k_1=0,1,\ldots, N_1-1,\quad k_2=0,1,\ldots,N_2-1\tag{2}$$
Plugging (2) into (1) gives
$$X[k_1+N_1k_2]=\sum_{n_2=0}^{N_2-1}\sum_{n_1=0}^{N_1-1}x[n_1N_2+n_2]W_N^{(n_1N_2+n_2)(k_1+k_2N_1)}\tag{3}$$
The complex exponentials in (3) can be simplified as follows:
$$W_N^{(n_1N_2+n_2)(k_1+k_2N_1)}=W_N^{N_2n_1k_1}W_N^{N_1N_2n_1k_2}W_N^{n_2k_1}W_N^{N_1n_2k_2}=\\= W_{N_1}^{n_1k_1}\cdot 1\cdot W_{N}^{n_2k_1}W_{N_2}^{n_2k_2}\tag{4}$$
because $W_N^{N_1}=W_{N_2}$, $W_N^{N_2}=W_{N_1}$, and $W_N^{N_1N_2}=W_N^N=1$. Plugging (4) back into (3) and rearranging terms gives
$$X[k_1+N_1k_2]=\sum_{n_2=0}^{N_2-1}\left[W_N^{n_2k_1}\left(\sum_{n_1=0}^{N_1-1}x[n_1N_2+n_2]W_{N_1}^{n_1k_1}\right)\right]W_{N_2}^{n_2k_2}\tag{5}$$
Equation (5) can be interpreted as follows: the result can be computed by computing $N_2$ length-$N_1$ DFTs of subsampled versions of the input sequence $x[n]$ (this is the inner sum in (5)), then multiplying the result by twiddle factors $W_N^{n_2k_1}$, and finally computing $N_1$ length-$N_2$ DFTs (the outer sum in (5)). This procedure is most easily visualized using matrices:
1. read the data row-wise into a $N_1\times N_2$ matrix
2. compute DFTs over all $N_2$ columns (this can be done in-place)
3. multiply the result element-wise with twiddle factors
4. compute DFTs over all $N_1$ rows (in-place)
5. read out the data column-wise to get the final length-$N$ result
This little octave/Matlab script illustrates the procedure:
N = 1536; % non-power of 2 FFT length used in LTE
N1 = 512; N2 = 3; % 512 x 3 = 1536
x = randn(N,1); % some input signal
% generate N1 x N2 matrix of input data, reading them row wise
X = reshape(x,[N2,N1]).';
% DFT over each column
X = fft(X);
% Multiplication with twiddle factors
X = X.*exp(-1i*2*pi/N*(0:N1-1)'*(0:N2-1));
% DFT over each row
X = (fft(X.')).';
X = X(:);
Matt's answer is excellent. For anyone wanting to do the same in Python, I thought I'd add the code (with a slight modification to use a sine wave input, which I prefer for debugging):
import numpy as np
N=1536
N1=512
N2=3 # 512 x 3 = 1536
nCycles = 2 # Number of cycles in source frame
time=np.arange(0,N ) # some input signal
x=np.sin(nCycles * time * 2. * np.pi / N)
# generate N1 x N2 matrix of input data, reading them row wise
x=np.reshape(x,(N1,N2))
# DFT over each column
X=np.fft.fft(x,axis=0)
# Element-wise multiplication with twiddle factors
X=np.multiply(X,np.exp(- 1j * 2. * np.pi / N * (np.arange(0,N1).T.reshape(N1,1) @ np.arange(0,N2).reshape(1,N2))))
# DFT over each row
X=np.fft.fft(X,axis=1) | 2022-08-09T14:53:07 | {
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https://math.stackexchange.com/questions/2487744/plot-a-complex-set-in-the-complex-plane/2487761 | # Plot a complex set in the complex plane
My lecturer only explained how to plot complex numbers on the complex plane, but he didn't explain how to plot a set of complex numbers.
I did some research online but I didn't find any clear explanation or method. I have an exercise to practice but I don't know how to even start!.
Any help would be really appreciated!
• Are you familiar with Analytic Geometry and the equations of circles, eclipses etc on the plane? – MathematicianByMistake Oct 24 '17 at 15:19
• Yes, how can I use them here? A friend explained me that with the equations(as in a.) I can obtain the radios and centre of a circunference. But I do not know how to do it, and the other cases I do not have any equation – Evoked Oct 24 '17 at 15:24
• Check the answer for a hint and a general approach. Your friend is correct. – MathematicianByMistake Oct 24 '17 at 15:25
HINT
Let's tackle the first one. A similar approach is required for the other two.
For $z=x+iy$, $x,y\in \mathbb{R}$ we get $$|z-1+i|=1\Rightarrow\\|x-1+i(y+1))|=1\Rightarrow\\\sqrt{(x-1)^2+(y+1)^2}=1\Rightarrow\\(x-1)^2+(y+1)^2=1$$
Does this-hopefully it does-remind you a more general equation of a circle?
Let's also look at the third one as well.
We have $$\operatorname{Re}\Big(\frac{z+1}{z-1}\Big)=\Re\Big(\frac{z+1}{z-1}\Big)=\Re\Big[\frac{(z+1)(\bar{z}-1)}{(z-1)(\bar{z}-1)}\Big]=\Re\Big[\frac{(z+1)(\bar{z}-1)}{|z-1|^2}\Big]\gt1$$ Now for $z=x+iy$, $x,y\in \mathbb{R}$ you can substitute on the last relationship and obtain an equation for $x,y$-an inequality to be precise.
• Yes it does! Since in the others there are inequalities, the representation would no longer be a circle? – Evoked Oct 24 '17 at 15:30
• I will help you get an equation for the last one as well, but-and don't take it personally-we can't do your homework for you, there would be no point in that if you don't try yourself. If the answer satisfies you you can upvote/accept.. – MathematicianByMistake Oct 24 '17 at 15:32
• This is not my homework! This is me trying to understand an exercise I searched for in a book! Please don't take in a bad way my words but I'm just trying to understand something, not asking for a particular class free – Evoked Oct 24 '17 at 18:32
• @Evoked No worries! I am only happy that the answer was helpful! Feel free to ask any questions on MSE. It is great for you that you wish to learn more and self-study. :-) – MathematicianByMistake Oct 24 '17 at 18:36 | 2019-06-26T08:17:05 | {
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https://mathhelpboards.com/threads/multidimensional-newton-raphson.27759/ | # Multidimensional Newton Raphson
#### Purplepixie
##### New member
Hello,
I am having difficulty in translating the univariate Newton's approximation {Xn = Xn-1 - [ f(Xn-1) / f'(Xn-1)]} into the multidimensional case. My multidimensional equation system is y = F.x where y and x are (nx1) column vectors and the coefficients matrix F is (nxn), so that (nx1) = (nxn).(nx1) = (nx1).
I have translated the univariate Newton's approximation to the n-variate case as:
x2 = x1 - (F.x1 / F'.x1) <=> x1 - (F.x1).Inv(F'.x1)
But F'.x1 is a nx1 column vector and has no inverse. I then thought that perhaps the x1's cancel out . But if so then we would have x2 = x1 - F.Inv(F') with the last term a nxn matrix, so (nx1) = (nx1) - (nxn), which is not possible.
Your assistance would be greatly appreciated!
(This is my first post incidentally, so pls excuse any breaches of protocol)
#### Klaas van Aarsen
##### MHB Seeker
Staff member
Hello,
I am having difficulty in translating the univariate Newton's approximation {Xn = Xn-1 - [ f(Xn-1) / f'(Xn-1)]} into the multidimensional case. My multidimensional equation system is y = F.x where y and x are (nx1) column vectors and the coefficients matrix F is (nxn), so that (nx1) = (nxn).(nx1) = (nx1).
I have translated the univariate Newton's approximation to the n-variate case as:
x2 = x1 - (F.x1 / F'.x1) <=> x1 - (F.x1).Inv(F'.x1)
But F'.x1 is a nx1 column vector and has no inverse. I then thought that perhaps the x1's cancel out . But if so then we would have x2 = x1 - F.Inv(F') with the last term a nxn matrix, so (nx1) = (nx1) - (nxn), which is not possible.
Your assistance would be greatly appreciated!
(This is my first post incidentally, so pls excuse any breaches of protocol)
Hi Purplepixie, welcome to MHB!
First we have to establish what $F'$ is.
If it is a matrix that does not depend on any variable, then $F'$ is the matrix that contains only zeroes.
Consequently we won't really get any useful result.
Newton's approach aims to solve $f(x)=0$, and makes use of the fact that $f(x) \approx f(x_0) + f'(x)(x-x_0) \implies x_0\approx x - \frac{1}{f'(x)}\cdot f(x)$.
We might try to solve $F(x_1,\ldots,x_n)=(0,\ldots, 0)$ in a similar fashion. In this case $F$ is not a matrix, but a set of $n$ functions. And each function has $n$ parameters.
If we take the derivative, we don't have just 1 derivative, but instead we have the derivatives of $n$ functions with respect to each of the $n$ parameters.
The result is a matrix of $n\times n$ derivatives. Let's call it $DF(x)$ to denote that it's a matrix of functions that depend on $x$.
Now we can write Newton's approximation as $x^{(k+1)} = x^{(k)} - \Big(DF(x^{(k)})\Big)^{-1}\cdot F(x^{(k)})$, where $x^{(k)}$ denote the successive approximations of $x$.
#### Purplepixie
##### New member
Dear Klaas,
Thank you for your clear and succinct explanation. I can now see where I was wrong. So much of mathematics would be simplified if representational systems could be improved!
All the best,
PP | 2020-09-30T03:53:59 | {
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"lm_q2_score": 0.8633916082162403,
"lm_q1q2_score": 0.8487874166153041
} |
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