A three-panel blanket set is described in which any two panels are releasably attachable to form a complete blanket. The first panel has a first insulation value and a first fastener section on its lower surface extending along the right side edge. The second panel has a second insulation value, a second fastener section on the upper surface extending along the left side edge, and a third fastener section on the lower surface extending along the right side edge. The third panel with a third insulation value has a fourth fastener section on the upper surface extending along the left side edge. The panel insulation values are different from each other. The first fastener section is releasably attachable to the second fastener section or to the fourth fastener section, and the fourth fastener section is releasably attachable to the first fastener section or to the third fastener section.

(1) Field of the Invention
The present invention relates to blankets with sections of different insulation properties, and in particular to blanket sets that include three panels configured so that selected pairs of the three panels can be joined to create a blanket with panels or sections with different insulation properties.
(2) Description of the Prior Art
When two occupants sleep in the same bed, it is often the case that the warmth requirement for one individual will be different from the other. That is, one individual may require a covering of a given insulation value, while the other individual may require a covering having a relatively lighter or heavier insulation value.
Various attempts to resolve this problem have been described in the prior art. Instead of using a blanket having a uniform thermal resistance or insulation value, the prior art discloses blankets comprised of first and second panels, each having a length equal to that of the desired complete blanket and a width approximately equal to one-half of the desired width of the blanket, with some means being provided to releasably attach the panels together at their inner edges to form the complete blanket. When used, the panel with the lesser insulation value is positioned over the side of the bed to be occupied by the person with the lesser warmth requirement, while the panel with the greater insulation value is positioned over the side of the bed to be occupied by the person with the greater warmth requirement. The adjacent edges of the two panels may be permanently attached by sewing, or releasably attached with a zipper, hook-and-loop fastener, button, ties, etc.
The prior art, i.e., U.S. Pat. No. 6,862,760 to Bradley et al., also broadly describes a three section comforter, i.e., two fabric sheets with insulation disposed between the sheets, that can be comprised of two selected sections that have releasable fasteners attached to both sides of each comforter section, with the fasteners preferably being identical and attached by an intermediate fastener, e.g., double buttons. However, Bradley et al. fail to describe a system specifically suitable for use with blankets. Therefore, there remains a need for a blanket set of three or more panels, constructed so that a selected pair of the panels can be easily and releasably attached by a zipper or other attachment means. There is also a need for a unique attachment means especially suited for use in attaching two blanket panels.


A current sensing circuit which provides for accurate in-line current sensing with extremely low insertion loss. A low valued resistor (e.g., 0.005 ohms) is connected in series with the source and load of the current to be measured. An analog-to-digital converter (ADC) is used to measure the resulting voltage generated across the resistor. In order to minimize inaccuracies due to voltage offsets introduced by the measurement circuitry when measuring the low voltages generated across such a low resistance, the ADC is "chopped," thereby causing self cancellation of any such offset voltages. A voltage source which provides a reference voltage for the ADC has a temperature coefficient which is approximately equal in magnitude to the temperature coefficient of the resistor. Hence, for a given current through the resistor, the reported voltage as measured across the resistor remains constant over temperature.

1. Field of the Invention
The present invention relates to current sensing circuits, and in particular, to current sensing circuits which provide for direct, in-line, non-inductive current sensing with low insertion loss.
2. Description of the Related Art
Unlike voltage monitoring, which can be done with a high degree of precision while introducing virtually no losses into the system whose voltage is being monitored, current sensing has required sacrificing precision for low insertion loss or suffering high insertion loss for greater precision. For example, one conventional current sensing technique with low insertion loss requires the use of an inductive sensor which senses current by the presence and strength of a magnetic field surrounding the conductor in which the current is conveyed. While the insertion loss introduced by this technique is very low (assuming negligible effects from interaction between the inductive sensor and the magnetic field generated by the current), a high degree of mechanical precision is required in that the proximity and placement of the inductive sensor must be consistent and carefully calibrated.
A conventional direct, in-line current sensing technique requires the insertion of a resistor in series with the load whose current is to be sensed. However, due to the necessity of having a relatively high voltage across such resistor so as to avoid errors caused by a variety of factors such as noise, such resistor can introduce significant losses due to its own power dissipation. While reducing the resistance will reduce such losses, the voltage generated by the current of interest is also reduced, thereby resulting in increased errors when such low voltages are measured. Such errors include inaccurate voltage readings due to offset voltages in the monitoring circuitry, as well as inaccuracies due to noise and temperature drift of the sensing resistor.
Accordingly, it would be desirable to have a current sensing circuit which allows high precision, in-line current sensing with low insertion loss.


A system has been developed for mapping the magnetic field during a balanced steady-state free precession (SSFP) sequence. Field maps are generated by analyzing the magnitude of images acquired using the SSFP imaging sequences with phase increment or frequency shift. The system maps the magnetic field relevant to the imaging sequence and need not rely on phase information. The field maps may be applied to adjust the hardware for correcting the field anomalies contained in the maps. The field maps may also be used for separation of water and fat signals in SSFP imaging.

Conventional magnetic resonance imaging (MRI) methods rely on a linear relationship between magnetic fields and spatial positions of the nuclear spins within the object being imaged. This relationship is realized by using a combination of homogeneous magnets and linear gradient coils. Deviations from this linear relationship can cause undesired image distortions and sometimes artifacts.
Field anomalies are deviations of magnetic fields from prescribed values. Anomalies can be static as well as sequence-induced. Static anomalies may occur from, for example, field inhomogeneities produced by the magnet used to generate the main static magnetic field. Sequence-induced anomalies may occur from, for example, eddy-currents or the Maxwell fields.
Two remedies often used for correcting anomalies in magnetic field are: (i) minimization of inhomogeneities by shimming the magnetic field using special shimming hardware, and (ii) post data acquisition image correction to compensate for the effects of the field anomalies. These remedies require accurate information about the spatial distribution of the magnetic field at the time when the image data is acquired.
Conventionally, maps of the magnetic field are generated using dedicated field-mapping sequences. Since images are acquired using a different sequence, such field maps contain no information about sequence-induced field anomalies relevant to the imaging sequence. In cases when the sequence-induced changes are significant, such as during balanced steady-state free precession (SSFP) sequences, it is preferable that the field maps be generated using the exact same sequence that is used for imaging.
In addition, field maps are often generated from the phases of the acquired images. As anyone skilled in the art understands, phases of acquired magnetic resonance (MR) signals can be affected by factors other than the static magnetic field. Such factors include, but not limited to, chemical shifts and radio-frequency (RF) field orientation. For example, in the single quadrature Dixon (SVD) method recently introduced for water-fat separated SSFP, maps of magnetic field are obtained from the phases of the isolated echo images. Since the phases of the echo signals are influenced by field inhomogeneities as well as chemical shifts, the magnetic field maps so generated are susceptible to interference from chemical shifts.
There is therefore a need for improved methods to generate magnetic field maps that are relevant to imaging sequences, wherein the maps are not degraded by interferences such as that from chemical shifts.

Magnetic field mapping during SSFP using phase-incremented or frequency-shifted magnitude images


Methods and user equipment (UE) apparatuses are provided to transmit a control signal using a first sequence or a second sequence in response to a detected control channel. The method includes identifying, through higher layer signaling, a configuration; determining a resource for control signal transmission from a sum of a lowest control channel element (CCE) index of the detected control channel and either a first offset or a second offset based on the identified configuration; and transmitting the control signal in the resource using the first sequence or the second sequence.

1. Field of the Invention
The present invention relates generally to wireless communication systems and more specifically, to transmission power control of sounding reference signals.
2. Description of the Art
A communication system includes a DownLink (DL) that conveys signals from at least one Transmission Point (TP) to User Equipments (UEs), and an UpLink (UL) that conveys signals from UEs to at least one Reception Point (RP). A UE, also referred to as a fixed or mobile terminal or a mobile station, includes a wireless device, a cellular phone, a personal computer device, and the like. A TP or an RP is generally a fixed station and is also referred to as a Base Transceiver System (BTS), a NodeB, an access point, and the like.
A communication system also supports several signal types of transmissions including data signals conveying information content, control signals enabling proper processing of data signals, and Reference Signals (RS), also known as pilots, enabling coherent demodulation of data or control signals or providing Channel State Information (CSI) corresponding to an estimate of a channel medium experienced by their transmissions.
DL data information is conveyed through a Physical DL Shared CHannel (PDSCH). DL Control Information (DCI) includes Scheduling Assignments (SAs) for Physical UL Shared CHannel (PUSCH) transmissions from UEs (UL SAs) or for PDSCH receptions by UEs (DL SAs). The SAs are conveyed through DCI formats transmitted in respective Physical DL Control CHannels (PDCCHs). In addition to SAs, PDCCHs may convey DCI that is common to all UEs or to a group of UEs. DCI also includes ACKnowledgment (ACK) information associated with a Hybrid Automatic Repeat reQuest (HARQ) ACK (HARQ-ACK) process transmitted to UEs from at least one TP through Physical HARQ-ACK Indicator CHannels (PHICHs) in response to respective receptions at RPs of data Transport Blocks (TBs) transmitted from the UEs.
UL data information is conveyed through a Physical UL Shared CHannel (PUSCH). UL Control Information (UCI) is conveyed through a Physical UL Control CHannel (PUCCH), unless a UE also transmits a PUSCH, in which case the UE may convey at least some UCI in the PUSCH. UCI includes HARQ-ACK information and is transmitted in response to a reception by a UE of data TBs. HARQ-ACK signaling is periodic or dynamic, if a respective reception of data TBs by a UE is semi-persistently (periodically) scheduled without a respective PDCCH or dynamically scheduled by a PDCCH. Other periodically transmitted UCI signaling includes DL CSI informing a NodeB of a channel medium experienced by a signal transmission to a UE and Scheduling Request (SR) informing a NodeB that a respective UE has data to transmit. A UL RS is used for demodulation of data or control signals, in which case the UL RS is referred to as DeModulation RS (DMRS), or to sound a UL channel medium and provide NodeBs with UL CSI, in which case it is referred to as a Sounding RS (SRS).
Typically, PDCCHs are a major part of a DL overhead. One method for reducing this overhead is to scale its size according to the resources required to transmit the PDCCHs and PHICHs in a DL Transmission Time Interval (TTI). Assuming Orthogonal Frequency Division Multiple Access (OFDMA) as the DL transmission method, a Control Format Indicator (CFI) parameter is transmitted through a Physical Control Format Indicator CHannel (PCFICH) to indicate a number of OFDM symbols allocated to a DL control region during a DL TTI.
FIG. 1 is a diagram illustrating a conventional structure for a DL control region in a DL TTI.
Referring to FIG. 1, a DL TTI includes one subframe having M symbols and a DL control region occupies a first N subframe symbols 110. The remaining M-N subframe symbols are primarily used to transmit PDSCHs 120. A PCFICH 130 is transmitted in some sub-carriers, also referred to as Resource Elements (REs) of the first symbol and conveys 2 bits indicating a PDCCH size of M=1, or M=2, or M=3 symbols. A PHICH 140 is also transmitted in some REs of the first symbol. Moreover, some symbols also contain RS REs, 150 and 160, that are common to all UEs for each of the TP antenna ports which in FIG. 1 are assumed to be two ports. The main purposes of UE-Common RS (CRS) are to enable a UE to obtain a channel estimate for its DL channel medium and to perform other measurements and functions. The remaining REs in the DL control region are used to transmit PDCCHs.
PDCCHs conveying SAs are not transmitted at predetermined locations in a DL control region and, as a consequence, each UE needs to perform multiple decoding operations to determine whether it has an SA in a DL subframe. To assist a UE with the multiple decoding operations, REs carrying a PDCCH are grouped into Control Channel Elements (CCEs) in the logical domain. For a given number of DCI format bits, a number of PDCCH CCEs depends on a channel coding rate, assuming Quadrature Phase Shift Keying (QPSK) as the modulation scheme. For UEs experiencing low or high DL Signal-to-Interference and Noise Ratio (SINR), TPs may respectively use a low or high channel coding rate for a PDCCH transmission in order to achieve a desired BLock Error Rate (BLER). Therefore, a PDCCH transmission to a UE experiencing low DL SINR may require more CCEs than a PDCCH transmission to a UE experiencing high DL SINR (different power boosting of REs of a CCE may also apply). Typical CCE aggregation levels for a PDCCH are, for example, of 1, 2, 4, and 8 CCEs.
For a PDCCH decoding process a UE determines a search space for candidate PDCCHs according to a common set of CCEs for all UEs (Common Search Space or CSS) and according to a UE-dedicated set of CCEs (UE-Dedicated Search Space or UE-DSS). The CSS may consist of the first NCCEUE-CSS CCEs in the logical domain. The UE-DSS is determined according to a pseudo-random function having as inputs UE-common parameters, such as a subframe number or a total number of CCEs in a subframe, and UE-specific parameters such as a Radio Network Temporary Identifier (RNTI). For example, for CCE aggregation levels Lε{1,2,4,8}, the CCEs for PDCCH candidate m are given by L·{(Yk+m)mod └NCCE,k/L┘}+i NCCE,k is a total number of CCEs in subframe k, i=0, . . . , L−1, m=0, . . . , M(L)−1, M(L) is a number of PDCCH candidates to monitor in a search space, and └ ┘ is the "floor" function rounding a number to its immediately smaller integer. Exemplary values of M(L) for Lε{1,2,4,8} are, respectively, {0, 0, 4, 2} in the UE-CSS, and {6, 6, 2, 2} in the UE-DSS. For the CSS, Yk=0. For the UE-DSS, Yk=(A·Yk−1)mod D where Y−1=RNTI≠0, A=39827 and D=65537.
PDCCHs conveying information to multiple UEs, such as a PDCCH conveying Transmission Power Control (TPC) commands for UEs for adjusting respective PUSCH or PUCCH transmission powers, are transmitted in the CSS. If enough CCEs remain in the CSS after transmitting PDCCHs conveying DCI to multiple UEs in a subframe, the CSS is also used to transmit PDCCHs providing SAs with specific DCI formats. The UE-DSS is exclusively used to transmit PDCCHs providing SAs. For example, the CSS may consist of 16 CCEs and support 2 PDCCHs with L=8 CCEs, or 4 PDCCHs with L=4 CCEs, or 1 PDCCH with L=8 CCEs and 2 PDCCHs with L=4 CCEs. The CCEs for the CSS are placed first in the logical domain (prior to an interleaving of CCEs).
FIG. 2 is a diagram illustrating a conventional PDCCH transmission process.
Referring to FIG. 2, after channel coding and rate matching, encoded bits of DCI formats are mapped to CCEs in the logical domain. The first 4 CCEs (L=4), CCE1 201, CCE2 202, CCE3 203, and CCE4 204 are used to transmit PDCCH to UE1. The next 2 CCEs (L=2), CCE5 211 and CCE6 212, are used to transmit PDCCH to UE2. The next 2 CCEs (L=2), CCE7 221 and CCE8 222, are used to transmit PDCCH to UE3. Finally, the last CCE (L=1), CCE9 231, is used to transmit PDCCH to UE4. The DCI format bits of a PDCCH is scrambled 240 with a binary scrambling code and are subsequently modulated 250. Each CCE is further divided into Resource Element Groups (REGs). For example, a CCE including 36 REs is divided into 9 REGs, each including 4 REs. Interleaving 260 is applied among REGs (blocks of 4 QPSK symbols). For example, a block interleaver is used with interleaving performed on symbol-quadruplets (4 QPSK symbols corresponding to 4 REs of a REG) instead of on individual bits. After REG interleaving, a resulting series of QPSK symbols is shifted by J symbols 270, and finally each QPSK symbol is mapped to an RE 280 in the DL control region of a subframe. Therefore, in addition to RS from TP antenna ports, 291 and 292, and other control channels such as a PCFICH or a PHICH 293, REs in a DL control contain QPSK symbols corresponding to DCI formats for UE1 294, UE2 295, UE3 296, and UE4 297.
After the reception of a PDSCH, a UE transmits HARQ-ACK information in a PUCCH to indicate the correct (ACK) or incorrect (NACK) reception of data TBs in a PDSCH.
FIG. 3 is a diagram illustrating a conventional structure for HARQ-ACK signal transmission in a PUCCH.
Referring to FIG. 3, HARQ-ACK signals, and RS enabling coherent demodulation of HARQ-ACK signals, are transmitted in one slot 310 of a PUCCH subframe including 2 slots. HARQ-ACK information bits 320 modulate 330 a Zadoff-Chu (ZC) sequence 340, for example using BPSK or QPSK, which is then transmitted after performing an Inverse Fast Fourier Transform (IFFT) operation. Each RS 350 is transmitted through an unmodulated ZC sequence.
For a UL system BandWidth (BW) including NRBmax,UL Resource Blocks (RBs), where each RB includes NscRB=12 REs, a ZC sequence ru,v(α)(n) is defined by a Cyclic Shift (CS) α of a base ZC sequence ru,v(n) according to ru,v(α)(n)=ejαnru,v(n), 0≦n<MscRS, where MscRS=mNscRB is a length of a ZC sequence, 1≦m≦NRBmax,UL, and ru,v(n)=xq (n mod NZCRRS) where the qth root ZC sequence is defined by
                    x        q                          (        m        )              =          exp                    (                                            -              jπ                                                                                qm                                            (                                  m                  +                  1                                )                                                          N            ZC            RS                          )              ,0≦m≦NZCRS−1 with q given by q=└q+½┘+v·(−1)└q┘ and q given by q=NZCRS·(u+1)/31. The length NZCRS of a ZC sequence is given by the largest prime number such that NZCRS<MscRS. Multiple RS sequences are defined from a single base sequence through different values of α. A PUCCH transmission is assumed to be in one RB (MscRS=NscRB).
FIG. 4 is a diagram illustrating a conventional transmitter for a ZC sequence.
Referring to FIG. 4, a mapper 420 maps a ZC sequence 410 to REs of an assigned transmission BW as they are indicated by RE selection unit 425. Subsequently, an IFFT is performed by IFFT unit 430, a CS is applied to the output by CS unit 440, followed by scrambling with a cell-specific sequence using scrambler 450, a Cyclic Prefix (CP) is inserted by CP insertion unit 460, and the resulting signal is filtered by filter 470. Finally, a transmission power PPUCCH is applied by power amplifier 480 and the ZC sequence is transmitted 490.
Different CSs of a ZC sequence provide orthogonal ZC sequences. Therefore, different CSs of a same ZC sequence are allocated to different UEs in a same PUCCH RB and achieve orthogonal multiplexing for respective HARQ-ACK signal and RS transmissions. Orthogonal multiplexing can also be in the time domain using Orthogonal Covering Codes (OCCs) where PUCCH symbols used for HARQ-ACK signal or RS transmission in each slot are respectively multiplied with a first OCC and a second OCC. For example, for the structure in FIG. 3, HARQ-ACK signal transmission is modulated by a length-4 OCC, such as a Walsh-Hadamard (WH) OCC, while RS transmission is modulated by a length-3 OCC, such as a Discrete Fourier Transform (DFT) OCC. In this manner, the multiplexing capacity is increased by a factor of 3 (determined by the OCC with the smaller length Noc). The WH OCCs, {W0, W1, W2, W2}, and DFT OCCs, {D0, D1, D2}, are:
            [                                                  W              0                                                                          W              1                                                                          W              2                                                                          W              3                                          ]        =          [                                    1                                1                                1                                1                                                1                                              -              1                                            1                                              -              1                                                            1                                1                                              -              1                                                          -              1                                                            1                                              -              1                                                          -              1                                            1                              ]        ,                      [                                                  D              0                                                                          D              1                                                                          D              2                                          ]        =                  [                                            1                                      1                                      1                                                          1                                                      ⅇ                                                      -                    j2π                                    /                  3                                                                                ⅇ                                                      -                    j4π                                    /                  3                                                                                1                                                      ⅇ                                                      -                    j4π                                    /                  3                                                                                ⅇ                                                      -                    j2π                                    /                  3                                                                    ]            .
Table 1 presents a mapping of a PUCCH resource nPUCCH to an OC noc and a CS α for a HARQ-ACK signal and RS transmission. For brevity, the RS associated with the HARQ-ACK signal will not be mentioned in the following. As a PUCCH is assumed to be transmitted over 1 RB including 12 REs, there is a total of 12 CS for a ZC sequence.
TABLE 1HARQ-ACK Signal and RS Resource Mapping to OC and CSOC noc for HARQ-ACK Signal and for RSCS αW0, D0W1, D1W3, D20nPUCCH = 0nPUCCH = 121nPUCCH = 62nPUCCH = 1nPUCCH = 133nPUCCH = 74nPUCCH = 2nPUCCH = 145nPUCCH = 86nPUCCH = 3nPUCCH = 157nPUCCH = 98nPUCCH = 4nPUCCH = 169nPUCCH = 1010nPUCCH = 5nPUCCH = 1711nPUCCH = 11
A UE can determine a conventional PUCCH resource nPUCCH for its HARQ-ACK signal transmission either through explicit signaling from serving TP(s) or through implicit signaling. The latter is based on CCEs used to transmit a PDCCH conveying a respective DL SA in response to which a UE transmits a HARQ-ACK signal. A one-to-one mapping may exist between conventional PUCCH resources used to transmit HARQ-ACK signals and CCEs used to transmit PDCCHs. For example, for UEs with one transmitter antenna port and a PDCCH transmission structure as in FIG. 2, a UE determines a conventional PUCCH resource for HARQ-ACK signaling from a CCE with a lowest index from a respective DL SA transmission. Then, UE1, UE2, UE3, and UE4 can respectively use PUCCH resource 1, 5, 7, and 9. If all PUCCH resources within an RB are used, resources in the immediately next RB are used. In general, a UE can determine a conventional PUCCH resource nPUCCH for HARQ-ACK signaling as nPUCCH=nCCE+NPUCCH where nCCE is a CCE with a lowest index for a respective DL SA and NPUCCH is TP-specific offset that is informed to UEs by higher layer signaling.
Improving coverage and cell-edge throughput are key objectives in communication systems. Coordinated Multi-Point transmission/reception (CoMP) is an important technique to achieve these objectives. CoMP operation relies on the fact that when a UE is in a cell-edge region, it is able to reliably receive a signal combined at a set of TPs (DL CoMP) and reliably transmit a signal combined at a set of RPs (UL CoMP). DL CoMP schemes can range from simple ones of interference avoidance, such as coordinated scheduling, to more complex ones requiring accurate and detailed channel information such as joint transmission from multiple TPs. UL CoMP schemes can also range from simple ones where PUSCH scheduling is performed considering a single RP to more complex ones where received signal characteristics and generated interference at multiple RPs are considered.
FIG. 5 is a diagram illustrating a conventional UL CoMP operation.
Referring to FIG. 5, a signal transmitted by a UE 510 is received from two RPs, RP1 520 and RP2 530. Scheduling coordination between the two RPs and combining of the respective received signals is facilitated by a fast backhaul link such as an optical fiber link.
Support of UL CoMP introduces new requirements for HARQ-ACK signaling in a PUCCH. As conventional HARQ-ACK signaling scrambles a respective ZC sequence with a respective RP-specific (cell-specific) sequence, it is not possible to support orthogonal multiplexing of HARQ-ACK signals in a same RB for reception at multiple RPs. For this reason, separate PUCCH RBs should be used for UL CoMP reception of HARQ-ACK signals. The scrambling of such HARQ-ACK signals is with a scrambling sequence that is common for all RPs constituting a set of UL CoMP RPs for a respective set of UEs (a CoMP-set specific ZC sequence is UE-specific and signaled to a UE by higher layer signaling).
The need to provide non-conventional PUCCH resources to support orthogonal multiplexing of HARQ-ACK signals for reception at multiple RPs, relative to conventional PUCCH resources supporting orthogonal multiplexing of HARQ-ACK signals for reception at a single RP, is associated with a respective increase in the UL overhead which reduces UL throughout.
Additionally, if both a conventional PDCCH and non-conventional PDCCH types are transmitted in a same subframe, collisions of PUCCH resources may occur, if the respective CCEs of the various PDCCH types are independently indexed and non-conventional PUCCH resources may then need to be configured for each of the non-conventional PDCCH types.
Therefore, there is a need to reduce an overhead resulting from assigning non-conventional PUCCH resources for transmissions of HARQ-ACK signals.
There is another need to provide mappings for compressing an amount of non-conventional PUCCH resources for HARQ-ACK signaling.
Finally, there is another need for indicating to a UE whether to use a conventional PUCCH resource or a non-conventional PUCCH resource for its HARQ-ACK signaling.

This patent is subject to a terminal disclaimer.

Transmission of acknowledgement signals from a user equipment for orthogonal reception at multiple points


A new and distinct cultivar of Ivy Geranium plant named ‘Dueparest’, characterized by its upright to outwardly spreading plant habit; freely basal branching habit; freely flowering habit; red and pink bi-colored double flowers; and good garden performance.

The present invention relates to a new and distinct cultivar of Ivy Geranium plant, botanically known as Pelargonium peltatum, and hereinafter referred to by the name 'Dueparest'.
The new Ivy Geranium plant is a naturally-occurring whole plant mutation of a proprietary selection of Pelargonium peltatum identified as code number E-16-18, not patented. The new Ivy Geranium plant was discovered and selected by the Inventor as a single flowering plant from within a population of plants of the parent selection in a controlled greenhouse environment in Rheinberg, Germany in May, 2007.
Asexual reproduction of the new Ivy Geranium plant by vegetative terminal cuttings in a controlled greenhouse environment in Rheinberg, Germany since May, 2007, has shown that the unique features of this new Ivy Geranium plant are stable and reproduced true to type in successive generations.


A drying apparatus for drying a semiconductor wafer includes a processing chamber including a rinsing section and a drying section adjacent thereto. The rinsing section has a chamber loading slot associated therewith for receiving the semiconductor wafer. The drying section has a chamber unloading slot associated therewith for outputting the semiconductor wafer. An exhaust control cap is carried by the processing chamber and includes a bottom wall, a top wall, at least one intermediate wall between the bottom and top walls, and a side wall coupled to the top, bottom and the at least one intermediate wall to define stacked exhaust sections. The exhaust control cap has a cap loading slot aligned with the chamber loading slot, a cap unloading slot aligned with the chamber unloading slot, and at least one exhaust port configured to be coupled to a vacuum source.

Chemical-mechanical polishing (CMP) is performed in the processing of semiconductor wafers. A standard CMP apparatus has a circular polishing pad and a rotating carrier for holding a semiconductor wafer. An abrasive slurry is used on the polishing pad. After a CMP operation, residual particles are left on the surface of the semiconductor wafer. These residual particles need to be removed.
Semiconductor wafers are typically cleaned in a cleaning apparatus which includes one or more brush stations each having a pair of rotary brushes for cleaning the major surfaces of the wafers. A conventional cleaning apparatus 10 is illustrated in FIG. 1. The cleaning apparatus 10 includes a pair of brushes 11 that clean the major surfaces of a semiconductor wafer 14 placed therebetween. The semiconductor wafer 14 is supported by rollers 16, which are also used to rotate the wafer. The cleaning apparatus 10 includes at least one spray bar 18 to direct a spray of fluid towards the semiconductor wafer 14.
As semiconductor processes continue to achieve smaller line widths to create semiconductor wafers with greater capacity, post CMP defects control becomes more important for improving wafer yield and reliability. A key factor in brush cleaning is to precisely control the distance between the brush and the semiconductor wafer. If the separation distance is too tight, the residual particles from the CMP operation will scratch 20 the semiconductor wafer 14 as they are brushed off, as illustrated in a highlighted section 21 of the wafer as provided in FIG. 2. If the separation distance is to loose, some of the residual particles will not be removed. For example, a residual particle 22 may cause a short between two lines 24 and 26, as illustrated in a highlighted section 27 of the wafer as also provided in FIG. 2.
After the cleaning apparatus, the semiconductor wafer is placed in a drying apparatus. A conventional drying apparatus is also referred to as a spin-rinse/dryer, and includes a processing chamber comprising a rinsing section and a drying section adjacent thereto. The sequence in the drying apparatus includes spinning the semiconductor wafer in the rinsing section while rinsing with a rinsing liquid, and then spinning the wafer in the drying section to dry. The rinsing section is to remove any residual particles not removed by the cleaning apparatus, whereas the drying section is to remove any residue marks left on a semiconductor wafer after rinsing.
Vapor within the drying section needs to be contained and exhausted so as to avoid contamination of the semiconductor wafer as it is removed from the drying apparatus. As the semiconductor wafer is removed, DI water or an IPA flow rate spike may cause a water mark or IPA to condense on the semiconductor wafer. This will lead to wafer yield and reliability being decreased. One approach to avoiding contamination is disclosed in U.S. Pat. No. 7,513,062. An exhaust control cap is provided for the drying apparatus and includes an exhaust manifold and a nitrogen blanket manifold. A gas flow module is coupled to the exhaust manifold and the nitrogen blanket manifold to control the vapor flow rate, the exhaust rate and the nitrogen blanket flow rate. In addition, an exhaust line may maintain a vertical laminar flow through the output portion, as well as diluting any vapor that may escape from the drying section.
Further developments to contain and exhaust vapor from the drying section of a drying apparatus are still desired. This is particularly so if the rinsing liquid within the rinsing section is toxic.

Drying apparatus with exhaust control cap for semiconductor wafers and associated methods


An automatic nailing head for an automatic nailing device is disclosed. A nailing head including a means for causing an angle between a central shank axis of the nail and a central axis of the throughbore to gradually decrease as the nail passes through a nail feed channel and into the throughbore is described. A picker plate capable of receiving a nail when in a first position and dropping the picked nail when moved to a second position is described. A nailing head with a quick release body portion for unjamming the heads without the use of tools is also described.

The present invention relates generally to automatic nailing devices and more particularly to a nailing head assembly for an automatic nailing device.
Automatic nailing devices, such as the device disclosed in the Colson U.S. Pat. No. 3,945,549, are known in the art. These devices advantageously replace manual labor in certain repetitive manufacturing operations such as in the manufacture of pallets and wooden fencing, for example. However, the automatic nailing unit disclosed in the Colson patent has several features that limit its capacity.
The automatic nailing device of Colson forms stringer pallets and is shown in FIG. 1. The device 20 includes a frame 22 supported by a pair of vertical support legs 24 at a feed end 26, and a pair of vertical support legs 28 at a discharge end 30.
Mounted to the frame 22 is a feed conveyor 32. Although only one side of the conveyor 32 is shown in FIG. 1, the opposite side is identical, except that the opposite side is a mirror image of the side shown. Only the side shown will be described for purposes of clarity.
The conveyor 32 includes a pair of spaced apart endless chains 34, mounted on sprockets 36, 38 and 40. The sprockets are mounted onto shafts (not shown) mounted for rotation in the frame by means known in the art. A drive motor (not shown) is provided to cause the chains to move along an endless path. A plurality of drive bars (not shown) are attached at each end to the chains, and are positioned transverse to a machine direction shown by arrow 42.
A flat plate 44 is mounted to an upper surface of the frame 22 and supports a plurality of stringer magazines 46 mounted above the plate 44. Each magazine 46 drops a stringer positioned in the machine direction 42 onto the drive bars (not shown). The drive bars advance a group of stringers in the machine direction 42 in response to instructions from a controller (not shown).
A plank magazine 48 is positioned to deposit planks in a direction transverse to the machine direction 42 and onto an upper surface of each stringer. The spacing between planks and the number of planks placed on the stringers is determined by the controller (not shown).
Next, the assembly of stringers and planks is moved under a nailing station 50. The nailing station 50 includes an upright frame 52 which supports a plurality of nailing units 54. The upright frame 52 is adjustable in height to accommodate different pallet thicknesses. One nailing unit 54 is provided for each stringer.
Each nailing unit includes a vertical plate 56 and a pneumatic cylinder 58 mounted vertically to the plate. A piston 60 of each cylinder 58 is positioned on a lower end of the cylinder 58 and is fixedly connected to a driving rod 62 by means of a coupling (not shown) an upper portion of the driving rod 62 includes an enlarged portion 64 for contacting an actuator 65.
Each nailing unit includes a nailing head 66 having a vertical throughbore for accepting the driving rod 62, and a nail feed channel fluidly connected to the vertical throughbore for depositing nails into the nailing head 66. Details of the known nailing head 66 will be described in more detail below.
A vibratory nail feeder 68 is mounted to the frame 22 and is provided for delivering vertically oriented nails to a nail track 70. The nail track delivers nails to the nail delivery channel of the nailing head 66.
A cross-sectional view of the nailing head 66 is shown in FIG. 2. A nail guide track 70 comprising a pair of declining, horizontally spaced bars 72 is positioned to deliver nails to the nailing head 66. The head 66 has a main body portion 74 having a vertically oriented slot 76 extending through a side wall for receiving nails from the guide track 70. The slot 76 defines an opening to a nail feed channel 78 which delivers nails by means of gravity to a vertical throughbore 80. The throughbore 80 is substantially cylindrical and intersects and is fluidly connected to the nail feed channel 78 through an elongated opening 82. The nail drops into the throughbore 80 and continues to fall until a leading edge contacts a pair of jaws 84a, 84b pivotally mounted to the body portion 74 by means of pivot pins 86a and 86b. The jaws 84a, 84b are spring biased by means of opposing leaf springs 88a and 88b.
A tensioning device 90 is provided to adjust the amount of force delivered to the jaws 84a, 84b by the leaf springs 88a, 88b. The tensioning device comprises a housing 92, a spring 94 and an adjustment knob 96 for adjusting the amount of force the spring 94 applies to the leaf springs 88a, 88b.
A driving rod 62 is driven downwardly by means of a hydraulic ram 60 (shown in FIG. 1) to force a nail which is retained in a lower portion of the throughbore 80 through the spring biased jaws 84a, 84b and into the materials to be nailed together (not shown) located beneath the jaws 84.
A nail picker plate 98 is mounted for horizontal, reciprocal sliding motion in the nailing head 66. The direction of motion of the picker plate is into and out of the paper. A top elevational view of the picker plate is shown in FIG. 3.
The picker plate has a first major surface 100, an opposite surface 102 (shown in FIG. 2) substantially flat edges 104, 106 and 108, and a cammed edge 110 for contacting a cammed surface on an actuator (not shown). The picker plate has a notch 112 for receiving a nail picked from the nail track 70. The picker plate also has an aperture 114 extending through the first and opposite major surfaces 100 and 102. Connecting the aperture 114 and the notch 112 is a channel 116 which permits the nail in the notch to move to the aperture. The notch is smaller in width than the nail head, and the nail is suspended in the slot 76 of the main body portion when positioned in the notch 112.
An actuator (not shown) causes the picker plate 98 to move in a direction shown by arrow 118 when a lower shoulder (not shown) on an enlarged portion 64 of the nail driving rod 62 contacts an upper surface of the actuator (not shown). A camming surface of the actuator which contacts camming surface 110 causes the picker plate to move in a direction shown by arrow 118 until a leading edge of the notch 112 is aligned with an open end of the nail track 70. The vibration of the nail feeder 68 and nail track 70 causes the nail to travel into the notch 112. This occurs as the driving rod 62 nears the lowest portion of the downstroke. As the driving rod 62 begins moving upwardly, the actuator releases pressure on the camming surface 110 and the picker plate 98 moves in a direction opposite the direction shown by arrow 118. During this portion of the movement, the picked nail moves through the channel, and into the aperture 114. The aperture is larger than the size of the nail head and the nail falls through the aperture into the nail feed channel 78 (shown in FIG. 2).
Referring back to FIG. 2, the nailing head 66 as well as the hydraulic cylinder 58 (shown in FIG. 1) are mounted onto the vertical plate 56. The vertical plate 56 is removably attached to the frame by means of an upper pair of horizontally positioned channel members (not shown) and a lower pair of horizontally positioned channel members 119a and 119b which are fixedly mounted to the vertical portions of the upright frame 52 (shown in FIG. 1).
The nailing head 66 is formed from a back section 120 having a vertical, cylindrical groove defining part of the inner surface of the throughbore 80. The head also has a lower front portion 122 having a vertically oriented cylindrical groove defining a portion of the throughbore and a portion of the nail feed channel 78. The head also has an upper front portion 124 having an inner throughbore defining a portion of the nail feed channel 78. Although these parts can be disassembled with the use of tools, it is time consuming and difficult to take the nailing head 66 apart to clear the throughbore 80 and nail feed channel 78 of nails plugging the passages. When the nail head is in a central portion of the nailing station, this task becomes particularly difficult, and causes production delays.
A disadvantage of the device shown in the Colson patent is that a nail is not positioned in the notch of the picker plate until the driving rod is nearing the bottom of the downstroke. Because the picker plate is only in a position to pick a nail for a small portion of the nail hammering shaft's stroke, the picker plate does not always receive a nail from the nail guide track 70. The presence of substances such as galvanized coatings on the nails that cause the nails to stick together may aggravate this problem. Accordingly, the operating speed of the known automatic nailing device is limited.
Another disadvantage is that the nails do not have sufficient time to fall from the picker plate to the bottom of the nail driver bore at higher production rates. The descent of the nail is slowed when the nail is making the transition from the nail feed channel to the throughbore, because the nail must change direction abruptly at the intersection of the nail feed channel and the throughbore. Accordingly, the operating speed of the known automatic nailing device is limited.


A hub apparatus for a bicycle wheel with improved aerodynamic mechanism is provided. In one embodiment, the apparatus includes a hub body having a parabolic shape. The hub body is disclosed approximately about a center plane of the bicycle wheel and has a plurality of retaining holes for retaining one end of bicycle spokes.

Improving aerodynamics of a bicycle is paramount to cycling sports. Bicycle wheel and hub systems have been designed and redesigned over the years to meet aerodynamic demand, along with many other parts of a bicycle, such as fork, crank systems, headsets, aero bars, etc. Also, the other features, such as light weight, minimum mechanical parts, etc., are also desirable to a person skilled in the art.

Hub apparatus for bicycle wheel and method


A method is disclosed for the preparation of silver halide tabular emulsion grains, containing at least 75% chloride, wherein at least 50% of the total projected area of all the grains is provided by said tabular grains, and wherein said tabular grains exhibit an average aspect ratio of at least 5:1, an average thickness not greater than 0.5 micron and an average diameter of at least 0.6 micron, by a process comprising the following steps:. (1) preparing a dispersion medium containing a gelatino-peptizer, and performing the precipitation of the silver halide grains under the following conditions: PA1 (2) adding to the dispersion medium, during and/or after completion of the precipitation, an amount of a spectral sensitizing dye preferably ranging from 10.sup.-5 mole to 5.10.sup.-3 mole per mole of silver halide at a pH lower than 4.0; in this way adsorbed crystal habit modifier is totally or almost totally desorbed from the tabular emulsion grains; PA1 (3) removing excessive inorganic salts by a procedure starting at a pH lower than 4.0, e.g. ultrafiltration, or flocculation by a polymeric agent or an inorganic salt followed by several washings and redispersion.. (a) the chloride ion concentration in the dispersion medium is lower than 0.15 molar at the start of the precipitation and remains lower than 0.15 molar during the whole precipitation, and pH is maintained between 5.0 and 9.0; PA2 (b) the precipitation is performed in the presence of a crystal habit modifyer according to general formula Ia or Ib: ##STR1## wherein Z represents the atoms necessary to form a fused on aromatic carbocyclic or heterocyclic, unsubstituted or substituted ring, e.g. substituted with alkyl, alkenyl, aryl, alkoxy, hydroxy, mercapto, carboxy, amino or halogen, and PA2 R is hydrogen or a substituent as defined for ring Z; n is 1 or 0, and PA2 Q represents carbon in which case n=1, or Q represents nitrogen, in which latter case n=0. PA2 This crystal habit modifier can be added a several different stages of the precipitation. It can be added as a whole in one portion or in several fractions. Preferably this crystal habit modifier is an adenine derivative and most preferably it is adenine itself.

Under conventional precipitation conditions silver halide emulsion grains rich in chloride are obtained showing a cubic morphology with (100) crystal faces. Special measures are required to modify this usual crystal habit and this is especially the case in the preparation of tabular grains rich in chloride containing (111) parallel major faces. Several publications reveal the use of so-called "growth modifiers" or "crystal habit modifiers". Klein and Moisar, in Berichte der Bunsengeselschaft 67 (4), p. 349-355, report inhibition of the growth rate of silver chloride by purine bases such as adenine. Claes et al., J. Photogr. Sci. Vol. 21 (1973), p. 39-50, showed that growth modifiers can be used to precipitate octahedral and rhombic dodecahedral silver chloride crystals and they attributed the crystal habit modification to variations in surface hydration caused by those additives; typical examples of these modifiers included adenine, thiourea, hypoxanthine, benzimidazole and benzothiazole derivatives. The mechanism of the growth modifying action of adenine was studied in detail by Szucs in J. Signal AM Vol. 6 (1978) No 5 p. 381-405. In Paper III-13 on the International Congress of Photographic Science Wyrsch reported a triple jet precipitation in the presence of ammonia and Cadmium ions. Depending on pAg and pH control, octahedral (111), rhombododecahedral (110) and cubic (100) crystal habits could be obtained.
In more recent patent applications tabular chloride rich emulsion grains are contemplated. So Wey U.S. Pat. No. 4,399,215 discloses tabular silver chloride emulsion grains with an aspect ratio of at least 8:1 and parallel (111) major crystal faces. Precipitation conditions include the use of ammonia. Rather thick tabular grains are obtained. Wey U.S. Pat. No. 4,414,306 discloses tabular silver chlorobromide grains with at least an annular region where the molar ratio of chloride to bromide ranges up to 2:3.
In U.S. Pat. No. 4,400,463 Maskasky describes the preparation of a new crystallographic form of tabular silver halide grains rich in chloride by performing the precipitation in the presence of a special peptizer having a thioether linkage and an aminoazaindene growth modifier. A preferred growth modifier is adenine. In this example the reaction vessel contained a rather high amount of chloride (0.5 molar) before the start of the precipitation and pH was adjusted to 3.0.
Maskasky U.S. Pat. No. 4,713,323 discloses the preparation of thin tabular grains (less than 0.35 micron) by a precipitation technique wherein at least a 0.5 molar concentration of chloride ion is present in the reaction vessel at the start and oxidized gelatin is used containing less than 30 micromoles of methionine per gram. In a preferred embodiment, illustrated by examples, a growth modifier e.g. an aminoazaindene like adenine is used. In this case the reaction vessel contained 0.5 molar chloride ions and pH was adjusted to 4.0.
Tufano U.S. Pat. No. 4,804,621 describes a process for preparing chloride rich tabular grains in the presence of aminoazapyridine growth modifiers represented by a general formula of which adenine and derivatives are excluded. A preferred compound is e.g. 4-aminopyrazoloE3,4,dlpyrimidine. In this example pCI was maintained at 0.7 and pH at 4.0. Houle and Tufano EP 0 430 196 discloses a method of stabilizing the crystal habit of this type of grains using outer bromide shells. However it can be expected that such a procedure will detoriate some of the specific advantages of chloride rich emulsions, e.g. the fast developability.
Unexamined Japanese Patent Publication (Kokai) 62-218959 describes a method for the preparation of chloride rich emulsion grains with average aspect ratio greater than 5:1 and (111) principal faces for at least 60% of the total grain projection area characterized by the presence of thiourea derivatives. The method is claimed to eliminate the stage of seed crystal preparation. In Unexamined Japanese Patent Publication (Kokai) 01-250943 a photographic material containing similarly prepared grains is claimed; the emulsion layer further contains a sensitizing dye.
In Takada U.S. Pat. No. 4,783,398 tabular grains rich in chloride are disclosed with an aspect ratio ranging from 2:1 to 10:1: they are preferably formed in the presence of a compound represented by a heterocyclic ring containing a sulphur atom.
The preparation of similar grains with an average aspect ratio of at least 5:1 in the presence of a gold compound, preferably combined with a Rhodium or Iridium compound, are described in the Unexamined Japanese Patent Publication (Kokai) 63-213836.
Unexamined Japanese Patent Publication (Kokai) 63-218938 combines tabular grains rich in chloride with the presence of a benzotriazole or mercaptotriazole derivative in the emulsion.
Unexamined Japanese Patent Publication (Kokai) 01-159646 claims the combination of tabular grains rich in chloride prepared in the presence of a crystal habit controlling agent, and a phenol type cyan colour coupler. The material is colour processed using a bleach-fixing bath.
Nishikawa U.S. Pat. No. 4,952,491 discloses a photographic material comprising a silver halide emulsion layer containing tabular grains rich in chloride showing an aspect ratio of at least 2:1, said tabular grains having been precipitated in the presence of a crystal habit controlling amount of a spectral sensitizing dye before and during nucleation and during precipitation.
In The Journal of Photographic Science, Vol. 36, p. 182, Endo and Okaji try to formulate an empirical rule for the modification of the crystal habit of silver chloride to form tabular grains. They conclude a.o. that especially 1,3-dimethyl urea and 1,3-diethyl urea are effective agents in order to obtain high aspect ratio tabular crystals with (111) faces.
In European Patent Application No. 90202792.9. filed 19.10.90, a method is disclosed for the preparation of silver halide tabular emulsion grains, containing at least 75% chloride, wherein at least 50% of the total projected area of all the grains is provided by said tabular grains, and wherein said tabular grains exhibit an average aspect ratio of at least 5:1, an average thickness not greater than 0.5 micron and an average diameter of at least 0.6 micron, by a process comprising the following steps:
preparing a dispersion medium containing a gelatino-peptizer, and a heterocyclic compound according to general formula (Ia) or (Ib), or one of their tautomeric forms, in a concentration ranging from 10.sup.-4 to 10.sup.-2 molar, said medium being adjusted to a pH between 5.0 and 9.0 and to a pCI between 1.0 and 2.0 by means of a chloride ions providing salt; formula (Ia) and (Ib) are: ##STR2## wherein Z represents the atoms necessary to form a fused on aromatic carbocyclic or heterocyclic, unsubstituted or substituted ring, e.g. substituted with alkyl, alkenyl, aryl, alkoxy, hydroxy, mercapto, carboxy, amino or halogen, and PA1 performing a silver halide precipitation comprising at least one double jet step by introducing in said dispersion medium at least one solution containing chloride ions and at least one solution containing silver ions in such a way that pCl is maintained between 1.0 and 2.0. pH is maintained between 5.0 an 9.0 and the concentration of compound (Ia) or (Ib) is maintained between 10.sup.-4 molar and 10.sup.-2 molar; PA1 removing excess of soluble salts by a wash technique performed at a pH value comprised between 4.0 and 9.0; this wash technique is preferably ultrafiltration. PA1 (1) preparing a dispersion medium containing a gelatino-peptizer, and performing the precipitation of the silver halide grains under the following conditions: PA1 (2) adding to the dispersion medium, during and/or after completion of the precipitation, an amount of a spectral sensitizing dye preferably ranging from 10.sup.-5 mole to 5.10.sup.-3 mole, and most preferably ranging from 3.10.sup.-5 mole to 2.10.sup.-3 mole per mole of silver halide at a pH lower than 4.0; in this way adsorbed crystal habit modifier is totally or almost totally desorbed from the tabular emulsion grains with maintenance of the tabular habit; PA1 (3) removing excessive inorganic salts by a procedure starting at a pH lower than 4.0. e.g. ultrafiltration, or flocculation by a polymeric agent or an inorganic salt followed by several washings and redispersion.
R is hydrogen or a substituent as defined for ring Z; n is 1 or 0, and PA2 Q represents carbon in which case n=1, or 0 represents nitrogen, in which latter case n=0. PA2 (a) the chloride ion concentration in the dispersion medium is lower than 0.15 molar at the start of the precipitation and remains lower than 0.15 molar during the whole precipitation, and pH is maintained between 5.0 and 9.0; PA2 (b) the precipitation is performed in the presence of a crystal habit modifyer according to general formula Ia or Ib: ##STR3## wherein Z represents the atoms necessary to form a fused on aromatic carbocyclic or heterocyclic, unsubstituted or substituted ring, e.g. substituted with alkyl, alkenyl, aryl, alkoxy, hydroxy, mercapto, carboxy, amino or halogen, and PA2 R is hydrogen or a substituent as defined for ring Z; n is 1 or 0, and PA2 Q represents carbon in which case n=1, or Q represents nitrogen, in which latter case n=0. PA2 This crystal habit modifier can be added at several different stages of the precipitation. It can be added as a whole in one portion or in several fractions. Preferably this crystal habit modifier is an adenine derivative and most preferably it is adenine itself.
In a preferred embodiment the heterocyclic compound according to general formula (Ia) or (Ib) is an adenine derivative, and in a still more preferred embodiment the heterocyclic compound is adenine itself.
However in the practical application of the method described in this European application serious problems arise concerning the developability of photographic materials containing tabular grains thus prepared. It was stated experimentally that a major part of the crystal habit modifier remains adsorbed at the surface of the tabular grains and that this phenomenon is the cause of the poor developability of those grains. As a result exposed and developed photographic materials containing such grains show inferior sensitometric properties such as an insufficient sensitivity, gradation and maximal density.
The present invention provides a remedy to the deficiencies of the application of the teachings of the European patent application explained above.
It is an object of the present invention to provide an improved method for the preparation of tabular grains rich in chloride.
It is a further object of the present invention to provide photographic materials containing tabular grains rich in chloride, which are readily developable and show excellent sensitometric characteristics.

Method for the preparation of tabular emulsion grains rich in chloride


In a shifting device comprising a shiftable double member arrangement and a shifting element for shifting the double member arrangement, the shifting element is arranged at least partially at a side of and next to, the double member arrangement.

The invention resides in a shifting device comprising a shiftable double member arrangement and a shifting element for shifting the double member arrangement.
DE 100 37 401 A1 discloses a shifting device with a double gear structure which can be shifted by a first shifting unit. The first shifting unit includes a slide sleeve which is disposed between two rotatable gears forming the double gear arrangement. Furthermore, the shifting device includes a second shifting unit which is partially integral with the first shifting unit and comprises a speed synchronizing unit, by way of which each rotatable gear of the double gear arrangement can be coupled for rotation with the shaft on which the rotatable gears are supported. The speed synchronizing unit includes a speed synchronizing body disposed between the rotatable gears of the double gear arrangement and a speed synchronizing body arranged beside the double gear arrangement.
It is the object of the present invention to provide a shifting device which is shorter than conventional shifting devices of the type with which the present invention is concerned.


In a seesaw type electromagnetic relay comprising an armature carrying out a seesaw motion by excitation of an electromagnet, and a movable spring which carries out the seesaw motion in synchronism with the seesaw motion of the armature and which comprises a par of spring sections, at least one of the spring sections forks into two branches to form a pair of spring pieces one of which is bent downward. A pair of movable contacts is mounted at both free end parts of the movable spring on a bottom surface of the movable spring. One of the movable contacts is mounted at the free end part in the spring pieces on bottom surfaces thereof to form a pair of movable contact elements acting as a twin movable contact. A pair of fixed contacts faces the pair of movable contacts. In an open state, a distance between one of the movable contact elements in the twin movable contact and the fixed contact opposed to the twin movable contact is different from another distance between another of the movable contact elements in the twin movable contact and the fixed contact opposed to the twin movable contact. The fixed contact opposed to the twin movable contact may consist of a twin fixed contact having a difference in level.

This invention relates to an electromagnetic relay and, more particularly, to a seesaw type electromagnetic relay serving as a continuous contact.
A seesaw type electromagnetic relay of the type is disclosed as prior art in Japanese Unexamined Patent Prepublication of Tokkai No. Hei 3-222,230, namely, 222,230/1991. According to Tokkai No. Hei 3 222,230, the seesaw type electromagnetic relay is symmetrical in shape for an armature. The armature comprises a central part rotatably supported on a supporting point and a pair of armature arm parts which extends from the central part in opposite directions. The armature makes a seesaw motion by excitation of an electromagnet. The armature is supported by a supporting member at the central part. The seesaw type electromagnetic relay comprises a pair of movable springs which extend to a direction in parallel with the armature with spaces left therebetween. Each movable spring comprises a middle part fixed to the supporting member and a pair of spring sections extending in parallel with both armature arm parts. The middle part of the movable spring is always in electrically contact with a common contact.
Each spring section forks into two branches to form a pair of spring pieces. The spring pieces have two free end parts on which two movable contact elements are mounted. The two movable contact elements are collectively called a twin movable contact. Each twin movable contacts faces a fixed contact. For each movable spring, a set of one twin movable contact and one fixed contact serves as a make contact while another set of another twin movable contact and another fixed contact acts as a break contact.
The seesaw type electromagnetic relay is operable as either a transfer contact or a continuous contact. The transfer contact is a contact where the break contact opens before the make contact closes. The continuous contact is a contact where the make contact closes before the break contact opens. The continuous contact is also called a make-before-break contact.
In a case of making the conventional seesaw type electromagnetic relay operable as the continuous contact, both spring sections in each movable spring are bent downward. As a result, the seesaw type electromagnetic relay has stronger maximum load force and it results in requiring large load energy upon driving of the continuous contact. As a result, there is little margin for a matching between magnetic attraction force of the electromagnet and urging force of an urging spring. In addition, the conventional seesaw type electromagnetic relay serving as the continuous relay consumes large power upon driving of it.
Various other electromagnetic relays have already proposed in, for example, Japanese Unexamined Utility Model Prepublication or Jikkai Nos. Sho 60-168,243, Hei 2-18,246, and Hei 4-78,720. Those Prepublications disclose the electromagnetic relays where a distance between an movable contact element of a twin movable contact and a fixed contact is different from a distance between another movable contact element of the twin movable contact and the fixed contact. However, all of the Prepublications only disclose the electromagnetic relays each of which acts as the make contact.

Seesaw type electromagnetic relay serving as a continuous contact with a low power consumption


A floating body germanium (Ge) phototransistor and associated fabrication process are presented. The method includes: providing a silicon (Si) substrate; selectively forming an insulator layer overlying the Si substrate; forming an epitaxial Ge layer overlying the insulator layer using a liquid phase epitaxy (LPE) process; forming a channel region in the Ge layer; forming a gate dielectric, gate electrode, and gate spacers overlying the channel region; and, forming source/drain regions in the Ge layer. The LPE process involves encapsulating the Ge with materials having a melting temperature greater than a first temperature, and melting the Ge using a temperature lower than the first temperature. The LPE process includes: forming a dielectric layer overlying deposited Ge; melting the Ge; and, in response to cooling the Ge, laterally propagating an epitaxial growth front into the Ge from an underlying Si substrate surface.

1. Field of the Invention
This invention generally relates to integrated circuit (IC) fabrication and, more particularly, to a germanium (Ge) phototransistor and associated fabrication process that capitalizes on the floating body effect of a device insulated from an underlying substrate.
2. Description of the Related Art
In a co-pending application entitled METHOD OF EPITAXIAL GROWTH OF GERMANIUM PHOTODETECTOR FOR CMOS IMAGER, Lee et al., Ser. No. 11/069,424, filed Feb. 28, 2005, a method for growing single crystalline Ge film by liquid phase epitaxy is presented, along with a means of fabricating a PIN photo detector for infrared photon detection. However, it is noted that the interface between Ge single crystal and the bottom dielectric layer (Si3N4) is not perfect. This interface may potentially cause leakage current in the diode. This interface is likely the result of the Ge PIN diode using the entire Ge film, so that the bottom interface deteriorates the electrical performance.
It would be advantageous if an improved interface could be created between the epitaxial layer of a Ge active device and the underlying substrate insulator.
It would be advantage if a Ge device could be fabricated so as to minimize the effect of the substrate interface.
It would be advantageous if the epitaxial layer of a Ge active device could be more effectively insulated from an underlying Si substrate.

Germanium phototransistor with floating body


The technical features of the present invention are for a method for quickly and easily constructing a building made of soil such as a house or a retaining wall by means of a process of installing a temporary mold in a frame member made of a mesh, compacting pugging clay into the temporary mold in a bottom-up method by a mechanical extruder, and immediately removing the temporary mold and drying the pugging clay.

It is probably safe to say that the natural building materials such as soil, wood and stone have been used in constructing the buildings for a long time along with human history.
In particular, comparing to stone or wood, soil has been used as a building material since ancient times because it is a common material which can be obtained everywhere. However, because it is poor in the mechanical strength and durability as a building material and because it has a poor image in its appearance, it has not been used as a mainstream material as compared to stone or wood.
However, since various bio-friendly and eco-friendly characteristics of soil have been newly recognized, soil is receiving the attention as a building material in the future.
In other words, the processing materials such as concrete, which are widely used as building materials today, destroy the nature from the manufacturing process and consume a large amount of energy, which result in increase of the emissions of carbon dioxide, which is a pollutant, as well as concrete wastes generated during the dismantling process contaminate the soil and damage the natural ecosystems in the end. In comparison, even if the soil buildings are dismantled after the life of the building has ended, the soil returns to the nature as it is, which can emphasize that there is no risk of corrupting or contaminating the natural ecosystems. In this point, the soil can be said as a representative of environmentally friendly materials more than anything.
In addition, as all living organisms are born from the soil and return to the soil, all organisms are inseparable from the soil. Therefore, the building which is mainly composed of the soil function as a space of a comfortable life, which is based on the various characteristics of the soil good for the human body, such as the functions of automatic temperature and humidity control, sterilization and deodorization, and radiation of far-infrared, which we can say that the soil is a bio-friendly material.
As stated above, while the soil has various eco-friendly and bio-friendly advantages, due to plain and simple appearance as well as other issues in robustness and durability, the soil has been used as the building materials with improved strength and durability such as bricks by burning the soil or by processing the soil with settlement agents.
However, in these building materials, the natural characteristics of soil are weakened or changed, let's leave them out of the discussion. If we examine the conventional construction methods using the soil as a main material, it is roughly classified into the wet construction methods and the dry construction methods.
First, the dry construction method can be divided into masonry type and compaction type. The masonry type is a method of construction to lay up block bodies, that are made by molding the soil inserted into the frame, using the methods such as bricklaying while the compaction type is a method of construction where the soil is filled into a form and compacted using a pounder and others, and this process is repeated to fill and compact the soil for constructing the wall.
As explained above, as the dry construction method does not use the water during the construction process, any separated drying process is not needed, as well as it does not demand any special construction technology. For this reason, it has an advantage of easy-to-construct, but this type of structure is weak against the lateral forces such as shear force, which has inefficient aspects of unnecessarily requiring a thicker wall and extremely low seismic performance and this structural weaknesses prevents this construction method from application to the large size buildings higher than the second floor.
Next, the wet construction methods include the pugging clay paint-over method, soil spray method, and recently developed pugging clay filling method.
Of these methods, the pugging clay paint-over method that has been used since ancient times, is a method to construct the wall by repeating the manual process daubing with the pugging clay inside and outside the core material which was made by weaving the plants such as stalks of millet.
It can be said that this method has been progressed one step in view of having reinforced the shear force and the earthquake resistance which are the greatest weaknesses of the construction of soil by enhancing the core material, however, due to a matter of primitiveness and effort of manual labor, the soil spray method where lath steel meshes instead of the core material of plants are arranged on the form materials and the soil and water or pugging clay are sprayed for paint-over using the mechanical equipment such as in shotcrete is partially used these days.
However, the so called dry construction method using this method where spraying the soil and water separately, has a problem that the soil is washed away during the process of spray or the soil is not properly kneaded, while the so called wet construction method where spraying pugging clay, has a problem that the mechanical equipment has often troubles by stuck clay due to its viscosity, which limits its application only to covering the soil wall thinly on already constructed wall and cannot be utilized for regular wall construction.
Finally, the pugging clay filling method is a construction method which has been recently developed and registered by this inventor as the Korean Patent No. 10-1003371 (Dec. 16, 2010) as a prior art considering the problems in the conventional wet construction as described above.
The prior invention described above regards to the frame member (10) composed of cage shaped wire mesh which is connected by a number of grid nets (13) between the front nets (11) and the rear nets (12) which are installed oppositely each other as shown in FIG. 1. This frame member performs the functions of the form and the reinforced rebar at the same time, if compared to concrete work.
The above described frame member is able to function as a shape retaining member without the formwork, which results from the clay specific-viscosity and plasticity which are realized when the soil becomes pugging.
In other words, unlike the concrete kneading which has a distinguished liquidity, if the pugging clay is poured into the frame member due to its strong viscosity and plasticity compared to liquidity, the pugging clay not only sticks on dense nets composed of wire mesh but also is not easily pushed out of wire mesh due to the plasticity that tends to maintain its form as long as the external force does not act.
Therefore, it can be considered that the prior invention has devised a new construction method for innovative structural form making the most of properties of the clay, however, the clay-specific viscosity and plasticity act as a double edged sword in the actual construction process as described above, and result in the reduced efficiency, which should be improved.
In other words, if we try to pour the pugging clay into the frame member, the pugging clay gets stuck immediately to wire mesh by the viscosity as described above, and if we try to mince the pugging clay using a pestle, the shape gets only distorted due to its unique plasticity, which makes it difficult to fill the frame member densely.
After all, this method has a problem that the construction efficiency gets significantly reduced in terms of construction period and manpower, because the work must be carried out manually one by one with hand tools in order to pour properly the pugging clay into the frame member. The existing mechanical spray method also exposed the same problem in the conventional spray method.

Method for constructing building made of dried soil and temporary frame used in same


A process for reducing the sulfur and ash content of coal. Particulate coal is introduced into a closed heated reaction chamber having an inert atmosphere to which is added 50 mole percent NaOH and 50 mole percent KOH moist caustic having a water content in the range of from about 15% by weight to about 35% by weight and in a caustic to coal weight ratio of about 5 to 1. The coal and moist caustic are kept at a temperature of about 300.degree. C. Then, water is added to the coal and caustic mixture to form an aqueous slurry, which is washed with water to remove caustic from the coal and to produce an aqueous caustic solution. Water is evaporated from the aqueous caustic solution until the water is in the range of from about 15% by weight to about 35% by weight and is reintroduced to the closed reaction chamber. Sufficient acid is added to the washed coal slurry to neutralize any remaining caustic present on the coal, which is thereafter dried to produce desulfurized coal having not less than about 90% by weight of the sulfur present in the coal feed removed and having an ash content of less than about 2% by weight.

This invention relates to a method of removing organic and inorganic sulfur compounds from coal and other carbonaceous combustible materials.
The recent energy crisis has increased the consumption of coal in the United States. However, there are many problems which need to be solved concerning the use of coal, the most important of which is environmental pollution. Most coal found in the United States contains from 0.5 to 10 weight percent sulfur which when burned is emitted as sulfur dioxide causing serious pollution problems in the atmosphere. Moisture in the atmosphere combines with the sulfur dioxide to produce acid rain, the consequences of which are far reaching and deleterious.
Because of the adverse impact of large quantities of sulfur from burning coal, it is necessary to reduce substantially the amount of sulfur which is released to the atmosphere. Only a small fraction of the available coal can be burned directly without violating current pollution control regulations. Thus, methods are being developed either to decrease the amount of sulfur in the coal before it is combusted or to remove sulfur from flue gas. However, flue gas desulfurization is expensive because of the high cost of the capital equipment and the cost of maintaining that equipment. Precombustion processes include conventional physical cleaning such as by comminution, chemical treatment, magnetic separation and coal conversion such as gasification and liquification.
Most physical cleaning methods separate mineral impurities from coal, utilizing differences in density of coal and mineral matter. This method only removes coarse mineral particles which are easily released, while leaving the finer particles in the coal. Coal gasification and liquification are not yet fully developed and are expensive. Magnetic separation can remove only liberated particles but does remove some ash forming minerals in addition to some of the sulfur. Chemical cleaning methods, which remove both organic and inorganic sulfur, are more effective than physical cleaning methods and are generally more economical than gasification and/or liquification.
Sulfur in coal may be classified into two general types, organic and inorganic. Organic sulfur is chemically bounded to the coal hydrocarbon matrix can only be removed by chemical means. Inorganic sulfur is present in coal with pyrite (FeS.sub.2) and in small amounts as a sulfate (generally, calcium or ferrous sulfate). Pyritic sulfur may range from about 0.5 to abut 10 percent by weight with individual particles of pyrite ranging from a few microns to a few inches in diameter. Large, liberated pyrite particles can be removed by hand or by various mechanical cleaning methods, but the particles that are finely distributed in the coal matrix need first to be liberated by fine grinding. Grinding and overgrinding produces a large percentage of fine material which is expensive not only because grinding to fine materials is costly but also the finer the coal, the more difficult in processing thereafter.
A number of different processes have been developed for the removal of organic and inorganic sulfur from coal and other burnable carbonaceous materials for reducing the ash content thereof. One method is known as the Gravimelt Process. The method is based on treating one part of finely powdered coal with approximately ten parts of a fused alkali such as sodium hydroxide, potassium hydroxide or mixtures thereof at 300.degree. C. to 400.degree. C. for 20-70 minutes. After removal of the coal floating on top of the melt and washing extensively with water to remove residual alkali and reaction products, substantial reductions in the sulfur content can be achieved. Subsequent washing of the treated coal with dilute sulfuric acid removes much of the mineral matter and neutralizes any remaining alkali, leaving a product that is relatively low in sulfur and low in ash content. Approximately ten percent of the original sulfur is present and approximately two percent ash is present by weight of the final product.
In another process for removing sulfur, a slurry of finely divided coal is in a solvent of methylchloroform, carbon tetrachloride or tetrachloroethylene and 30-70 weight percent water is prepared. Gaseous chlorine is bubbled through the slurry at 60.degree.-130.degree. C. and from 0-60 psig for about 45-90 minutes to oxidize the coal. The process will remove about 60 percent of the total sulfur in the coal removing about 50 percent of the organic sulfur and about 70 percent of the pyritic sulfur. While the process is reasonably effective, it is not as effective as the Gravimelt Process leaving about 50 percent of the organic sulfur in the coal. In addition, substantial amounts of residual chlorine remain in the coal which can produce highly corrosive combustion products upon burning.
The patent to Aida et al. U.S. Pat. No. 4,497,636 issued Feb. 5, 1985 attempts to combine the Gravimelt technology and the chlorinating technology in which chlorine gas is used as an oxidant and thereafter the carbonaceous material separated from the liquid chlorine is contacted with molten caustic, the operating temperatures used being in the range of from about 250.degree.-400.degree. C. with about 325.degree. C. being preferred. The results reported are about 90 percent of the total sulfur present in the coal being removed. Lower concentrations of alkali have been used to remove sulfur as reported in the Reggel et al. U.S. Pat. No. 3,993,455 and alkali has been used in super critical fluid conditions supposedly to release sulfur as reported in the Narain et al. U.S. Pat. No. 4,775,387. The three patents referred to, these being the Reggel et al. patent, the Aida et al. patent and the Narain et al. patent are all assigned to the assignee of the present invention and illustrate a portion of the resources of the government devoted to finding better and cleaner methods for burning coal.

The United States of America as represented by the United States Department of Energy


The present application describes a retention device for anchoring a medical device within the vasculature. The device may include expandable member coupled to an intravascular medical device and proportioned for receipt within a vessel. At least a portion of the expandable member is expandable to radially engage a vessel wall and to thereby retain the medical device within the vessel. The system is suitable for a variety of intravascular devices, including but not limited to ICD's, pacemakers, and intravascular drug delivery systems.

Pacemakers, defibrillators and implanted cardioverter defibrillators ("ICDs") have been successfully implanted for years for treatment of heart rhythm conditions.
Pacemakers are implanted in patients who have bradycardia (slow heart rate). The pacemakers detect periods of bradycardia and deliver electrical stimuli to increase the heartbeat to an appropriate rate.
ICDs are implanted in patients who may suffer from episodes of fast and irregular heart rhythms called tachyarrhythmias. An ICD can cardiovert the heart by delivering electrical current directly to the heart to terminate an atrial or ventricular tachyarrhythmia, other than ventricular fibrillation. An ICD may alternatively defibrillate the heart in a patient who may suffer ventricular fibrillation (VF), a fast and irregular heart rhythm in the ventricles. During a VF episode, the heart quivers and can pump little or no blood to the body, potentially causing sudden death. An ICD implanted for correction of ventricular fibrillation will detect a VF episode and deliver an electrical shock to the heart to restore the heart's electrical coordination.
Another type of implantable defibrillation device treats patients who may suffer from atrial fibrillation (AF), which is a loss of electrical coordination in the heart's upper chambers (atria). During AF, blood in the atria may pool and clot, placing the patient at risk for stroke. An electrophysiological device implanted for correction of atrial fibrillation will detect an AF episode and deliver an electrical shock to the atria to restore electrical coordination.
Pacemakers and ICDs are routinely implanted in the pectoral region either under the skin (subcutaneous) or under the pectoral muscle. The leads are placed at appropriate locations within or on the heart. Because of this complexity, a cardiologist identifying a heart rhythm condition may be required to refer his or her patient to sub-specialists or surgeons for implantation of a pacemaker or ICD—thus delaying implantation of the device in a patient who urgently needs it. It is thus desirable to simplify these devices and the procedures for implanting them so as to permit their implantation by a broader range of physicians.
U.S. application Ser. No. 10/453,971, filed Jun. 4, 2003, and Ser. No. 10/862,113, filed Jun. 4, 2004, describe intravascular systems that may be used to deliver electrical energy to the heart such as for defibrillation, pacing, and/or cardioversion of the heart. These applications are incorporated herein by reference for all purposes.
Generally speaking, the systems described in the '971 and '113 applications include at least one housing containing the necessary circuitry and related components, and optionally include at least one lead carrying the electrodes needed to deliver the electrical energy to the body. Some or all of these components are positioned within the vasculature, such as in the superior vena cava ("SVC"), the inferior vena cava ("IVC"), the left or right subclavian vein ("LSV" or "RSV"), coronary sinus and/or within other vessels in the venous or arterial system. For some of the implant components (such as the housing and/or lead), retention devices are needed to retain the implant within the vasculature.
The '971 and '113 applications sleeves or anchors that may be introduced after the implant to be retained has been positioned at a desired location within the vessel, or that may be introduced simultaneously with the implant. In some embodiments, the anchor is positioned adjacent the implant and expanded to a radially expanded position. During this expansion step the retention sleeve may self-expand and/or it may be expanded using an expansion tool such as a balloon passed into the retention sleeve's central lumen and subsequently inflated. When the retention sleeve is expanded, its radial forces engage the implant and secure the implant against the vessel wall. Blood flowing through the vessel passes through the tubular interior of the retention sleeve.
The '971 and '113 applications describe the retention devices as being either separate components from the implants they are intended to retain, or as being integral with the retained components. The present application is directed to alternative retention devices for supporting medical implants within blood vessels. Although the focus of this description is on retention devices that are integrated with the medical implants they retain, they are equally useful as components that are physically separate or separable from the medical implants. It should also be noted that although these retention devices will be described in the context of intravascular cardioverters, defibrillators, and or pacemakers, they are equally suitable for any type of medical implant that must be retained within a blood vessel.

Method and apparatus for retaining medical implants within body vessels


An apparatus for preventing vacuum of a scroll compressor is provided. The apparatus includes a fixed scroll having a bypass hole, through which a suction chamber of a casing is connected to a discharge chamber of the casing; a valve block having an inner channel, through which the bypass hole of the fixed scroll is connected to the discharge chamber of the casing, and fixedly installed at the fixed scroll; and a valve member disposed between the bypass hole of the fixed scroll and the inner channel of the valve block. Since the apparatus is assembled at an outer periphery of the fixed scroll, a processing error of the fixed scroll is prevented, thereby reducing production costs. Also, since foreign materials generated when a refrigerant channel is processed do not remain in a valve hole, the valve member is prevented from being mal-operated, productivity is enhanced, and fabrication costs reduced.

In a scroll compressor, an orbiting scroll that performs an orbit movement by a crank shaft is provided with an orbit wrap of an involute shape thereby to form a compression chamber that consecutively moves with a fixed wrap of an involute shape of a fixed scroll. As the orbiting scroll performs an eccentric orbit movement by a rotation of the crank shaft, a volume of the compression chamber is decreased thereby to compress a refrigerant contained in the compression chamber. When the scroll compressor is operated with a certain compression ratio, the compressed refrigerant is discharged through an outlet.
The scroll compressor can implement a relatively higher compression ratio than any other compressor, and can generate a stable torque by smoothly performing a suction stroke, a compression stroke, and a discharge stroke of a refrigerant. Accordingly, the scroll compressor is being widely used to compress a refrigerant in an air conditioner, etc. When the scroll compressor is normally driven, a suction chamber of a low pressure and a discharge chamber of a high pressure are divided from each other. On the contrary, when the scroll compressor is abnormally driven in a state of a low suction pressure such as a pump down or a cycle clogging, a refrigerant of a high pressure discharged to the discharge chamber is introduced into a compression chamber through the suction chamber under a state that the suction chamber is connected to the discharge chamber, thereby preventing a high vacuum state of the compressor.
FIG. 1 is a sectional view showing an example of a scroll compressor having a vacuum preventing apparatus in accordance with the conventional art, FIG. 2 is a sectional view showing an operation state of the vacuum preventing apparatus when the scroll compressor is normally driven in accordance with the conventional art, and FIG. 3 is a sectional view showing an operation state of the vacuum preventing apparatus when the scroll compressor is driven at a high vacuum state in accordance with the conventional art.
As shown, the conventional scroll compressor comprises a casing 10 divided into a suction chamber S1 of a low pressure and a discharge chamber S2 of a high pressure; a main frame 20 fixedly installed in the casing 10; a fixed scroll 30 fixedly installed at an upper surface of the main frame 20 and having a compression chamber P that consecutively moves with an orbiting scroll 40; and a vacuum preventing unit 50 installed at the fixed scroll 30 for introducing a discharged refrigerant to the compression chamber P by connecting the discharge chamber S2 of a high pressure to the suction chamber S1 of a low pressure when a pressure of the compression chamber P is lowered into a pressure less than a certain pressure.
The casing 10 is divided into the suction chamber S1 of a low pressure and the discharge chamber S2 of a high pressure by a high-low pressure separation plate 11 of which an inner circumferential surface is adhered to an upper surface of a plate 31 of the fixed scroll 30 and an outer circumferential surface thereof is adhered to an inner circumferential surface of the casing 10. A suction pipe 12 is connected to the suction chamber S1, and a discharge pipe 13 is connected to the discharge chamber S2.
A bypass channel for opening and closing the suction chamber S1 and the discharge chamber S2 of the casing 10 by a sliding valve 51 of the vacuum preventing unit 50 is formed at the plate 31 of the fixed scroll 30. The bypass channel comprises a valve hole 32 concaved with a certain depth at the plate 31 of the fixed scroll 30 in a horizontal direcrion-4, a suction pressure channel 33 through which the valve hole 32 is connected to the suction chamber S2 of the casing 10, a middle pressure channel 34 through which the valve hole 32 is connected to the compression chamber, and a discharge pressure channel 35 through which the valve hole 32 is connected to the discharge chamber S2 of the casing 10.
The vacuum preventing unit 50 comprises a sliding valve 51 for closing the discharge pressure channel 35 by being slidably inserted into the valve hole 32 in a horizontal direction when the compression chamber P maintains a certain pressure, and for introducing a refrigerant inside the discharge chamber S2 into the suction chamber S1 by connecting the discharge pressure channel 35 to the suction pressure channel 33 when a pressure of the compression chamber P is lowered into a pressure less than the certain pressure; and a valve spring 52 provided at one side of the sliding valve 51 for elastically supporting the sliding valve 51.
Unexplained reference numeral 36 denotes a fixed wrap, 37 denotes an inlet, 38 denotes an outlet, 41 denotes an orbit wrap, 53 denotes a valve cover, 54 denotes a fixing pin, 61 denotes a stator, 62 denotes a rotor, 63 denotes a crank shaft, and 64 denotes a sub frame.
In the conventional vacuum preventing apparatus of a scroll compressor, when a new refrigerant is supplied to the compression chamber P, the vacuum preventing unit 50 closes between the discharge pressure channel 35 of the fixed scroll 30 and the suction pressure channel 33 and thus a refrigerant is normally compressed. On the contrary, when a little amount of refrigerant is sucked into the compression chamber P, the vacuum preventing unit 50 connects the discharge pressure channel 35 to the suction pressure channel 33 so that a refrigerant discharged to the discharge chamber S2 can be supplied to the compression chamber P via the suction chamber S1.
As shown in FIG. 2, when the compression chamber P maintains a certain pressure, the pressure of the compression chamber P becomes equivalent to a resultant force between an elastic force of the valve spring 52 and a pressure of the suction chamber S1, and thus the sliding valve 51 closes the discharge pressure channel 35. On the contrary, as shown in FIG. 3, when the pressure of the compression chamber P is lowered into a pressure less than the certain pressure, the pressure of the compression chamber P becomes smaller than the resultant force between an elastic force OT the valve spring 52 and a pressure of the suction chamber S1. As the result, the sliding valve 51 is moved towards an opposite side to the valve spring 52, and thus the discharge pressure channel 35 is connected to the suction pressure channel 33.
However, the conventional vacuum preventing apparatus of a scroll compressor has the following problems.
First, the valve hole 32 is concaved by a certain depth at an outer circumferential surface of the plate 31 of the fixed scroll 30 towards a center of the plate 31, and the middle pressure channel 34 or the discharge pressure channel 35 is connected to the valve hole 32. As the result, burr generated when the middle pressure channel 34 and the discharge pressure channel 35 are mechanically processed remains in the valve hole 32 thereby to serve as an obstacle when the sliding valve 51 is operated.
Second, oil and refrigerant inside the compression chamber P are introduced into the valve hole 32 through the middle pressure channel 34 connected to the valve hole 32, and then remain in the valve hole 32 thereby to serve as an obstacle when the sliding valve 51 is operated.
Third, since the valve hole 32 is formed at a lateral wall of the fixed scroll 30 in a horizontal direction, a processing of the valve hole 32 is difficult, a defective proportion in processing the fixed scroll 30 is increased, and thus a production cost for the fixed scroll 30 is increased. Furthermore, since a gap between the valve hole 32 and the sliding valve 51 is increased due to an inaccurate dimension of the valve hole 32, a refrigerant inside the discharge chamber S2 is leaked thereby to degrade an efficiency of the compressor.

Apparatus for preventing vacuum of scroll compressor