The negative electrode for lithium-ion secondary battery is used in which a product of tensile strength and thickness of a negative electrode having a negative electrode active material layer containing silicon and silicon oxide as main components is 3.8 to 9.0 N/mm and a value obtained by dividing the product of the tensile strength and the thickness of the negative electrode by a product of tensile strength and thickness of a negative electrode current collector is 1.06 to 1.29.

1. Field of the Invention
The present invention relates to a negative electrode for lithium-ion secondary battery and a lithium-ion secondary battery.
2. Description of the Related Art
Since lithium-ion secondary batteries are light and have a high capacity in comparison with nickel-cadmium batteries, nickel-hydride batteries, and the like, the lithium-ion secondary batteries are widely used as power sources for portable electronic devices. In addition, the lithium-ion secondary batteries are prime candidates as power sources mounted in hybrid vehicles and electric vehicles. In recent years, along with downsizing and high functionality of portable electronic devices, the lithium-ion secondary batteries with a higher capacity, which are the power sources of these portable electronic devices, are expected.
The capacity of the lithium-ion secondary batteries mainly depends on active materials of electrodes. As for negative electrode active materials, graphite is generally used. However, a theoretical capacity of graphite is 372 mAh/g, and batteries in practical use with a capacity of about 350 mAh/g are already used. Accordingly, in order to obtain nonaqueous electrolyte secondary batteries having sufficient capacity as energy sources of high-functional portable devices in the future, there are needs for achieving a higher capacity, and negative electrode materials having larger theoretical capacity than that of graphite are needed for achieving a higher capacity.
At present, alloy-based negative electrode materials such as silicon and silicon oxide have attracted attention. Silicon can electrochemically intercalate and deintercalate lithium ions and charging and discharging with a very large amount of capacity compared with graphite are possible. In particular, it is known that the theoretical discharge capacity of silicon is about 4210 mAh/g, which is 11 times the theoretical discharge capacity of graphite.
However, when silicon and silicon oxide are used as a negative electrode active material to achieve a high capacity, several problems arise. Specifically, when lithium is intercalated, a very large amount of volume expansion is accompanied due to a change from the original crystal structure by forming a lithium-silicon alloy. For this reason, exfoliation of a negative electrode active material layer or fracture of the negative electrode occurs, and there is a problem of remarkably large cycle deterioration.
In order to solve the above problem due to the expansion, the volume expansion is suppressed by using a negative electrode current collector having high tensile strength such that a 0.2% proof stress is 250 N/mm2 or more or tensile strength is 300 N/mm2 or more to suppress a fracture of the negative electrode, and adhesion between the negative electrode active material layer and the negative electrode current collector are improved by using the negative electrode current collector with a surface roughness Rz of 0.6 to 10 μm by a roughening treatment (for example, Japanese Unexamined Patent Application No. 2011-192563). Moreover, the thickness of the negative electrode active material layer is limited so that the tensile strength of the negative electrode current collector is 3.82 N/mm or more, and the tensile strength of the negative electrode current collector per 1 μm thickness of an active material thin film is 1.12 N/mm or more (for example, Japanese Patent 3733070).

Negative electrode for lithium-ion secondary battery and lithium-ion secondary battery


Methods and systems for the differential wear mapping of the surface of an object are described. The process involves establishing a radioactive wearing surface on the object through surface layer activation (SLA) and/or radioactive ion implantation (RII) and imaging the wearing surface with conformal photographic or x-ray film to produce a first autoradiograph providing a baseline profile of the wearing surface. After subjecting the object to wear during testing the wearing surface is again imaged with conformal photographic or x-ray film to produce a second autoradiograph providing a worn profile. The first and second autoradiographs of the wearing surface are digitized and compared to produce a differential wear map of the surface of the object.

1. Field of the Invention
The present invention relates generally to methods and systems for measuring wear on material surfaces. The present invention relates more specifically to methods and systems for measuring the wear on a surface utilizing contact autoradiographs of residual radiation established by surface layer activation.
2. Description of the Related Art
Efforts have been made in the past to provide systems and methods for tracking, measuring, and monitoring the degree of wear on a variety of material surfaces. Of particular interest is the degree of wear that occurs over time on surfaces that, by operation of the machinery they are associated with, encounter abrasive surfaces or other metallic surfaces with the result that over time a significant amount of wear can occur. Many such surfaces can be found on critical engine components and other elements of vehicle systems that repetitively, or periodically, move against one another in a manner that produces wear over time. Knowing the degree and pattern of wear in components like these can be critical to safety and maintenance protocols associated with the machinery and equipment the components serve. Such testing further promotes and facilitates the development of improved designs and manufacturing techniques for these mechanical structures.
A number of efforts to determine the degree of wear on a surface involve placing or depositing some measurable characteristic onto the surface that will gradually be worn away as the surface itself is worn away. In some of these systems, this deposition process involves the placement of an additional over-layer on the surface and other times it involves providing some measurable characteristic to a surface layer (or layers) of the material itself. One obvious drawback to the use of the former approach (depositing an additional layer) is the resultant change in the character of the surface over time as the deposition layer is worn away. A further drawback is the fact that this additional layer will likely change the wear characteristics of the material. Preference is given therefore to techniques that utilize the actual material that the surface is made up of already, and that instill some characteristic in the upper layers of the material which can be monitored and measured over time.
The simplest example of such a wear monitoring system would involve a material that can be modified in some visual form on the upper layer or layers of its surface. Often this takes the form of creating physical patterns in the surface layers of the material when it is manufactured such that over time the patterns are changed or become visibly exposed as the wear on the surface occurs. This procedure, however, also suffers from the fact that the surface thereby changes physical appearance and therefore its physical characteristics, which could ultimately have an effect on the manner in which the surface interacts with other surfaces that it comes in contact with.
It would therefore be preferable if a wear monitoring system did not result in a physically changed surface over time and maintained the physical surface characteristics (with the exception, of course, of the wear that occurs) and therefore does not affect the manner in which the part or component functions (again, at least not beyond the extent to which the wear itself affects the functionality of the component).
The process of instilling a material with a characteristic that resides only in a surface layer can take a number of different approaches or forms. It is possible for example to provide a color dye to a material, with coloration variation according to a depth away from the surface into the material. Some such components may for example incorporate red dyes deep within the material, with green or blue dyes residing within the upper surface layers of the material, or more simply the density of the dye may vary according to depth. As the surface is worn away, areas beneath the surface are exposed and, upon visual inspection, differences can be determined in the degree of wear over different parts of the component. Such visually inspectable systems have drawbacks in that the degree of wear generally has to be significant in order for distinctions to be made between the upper layers and lower layers of the material. In addition, many materials are simply not receptive of dyes; at least not to the extent that visual inspection might discern quantifiable information.
Autoradiography is a process whereby the radioactive properties of a material are mapped by imaging the surface of the material (typically with photographic or x-ray film, although any two-dimensional detector array will do) and measuring the radiation generated by each particular surface area component. Autoradiography refers to the process of measuring the radiation generated by the object itself (as opposed to reflected or absorbed radiation measurements) and is generally carried out in a manner that creates an image of the radiation source whereby the radiation intensity varies in different portions across the view being measured. Such systems are used, for example, to locate radiation leaks in closed systems, and other fields where low level radiation is generated by an object in a manner that requires monitoring or tracking.

Autoradiography-based differential wear mapping


Providing wireless communication between an access point and a wireless device includes the wireless device operating in a first reduced power mode for a first period of time, where the wireless device does not receive communications while operating in the first reduced power mode and includes the wireless device operating in a full power mode for a second period of time, where the wireless device transitions back into the first reduced power mode in response to there being no communication from the access point to the wireless device during the second period of time. The wireless device may be a wireless sensor. In response to the access point directing the wireless sensor to collect data at a particular collection time, the wireless sensor may enter the first reduced power mode for a time prior to the collection time.

1. Technical Field
This application relates to the field of electronic components, and more particularly to the field of communication and power management for wireless electronic components.
2. Description of Related Art
There are systems that use a plurality of sensors to collect data. For example, aircraft monitoring systems may use a plurality of sensors that provide sensed data to one or more central processors that use the data to assess the health of system components. See, for example, U.S. Pat. No. 6,847,917, which is incorporated by reference herein, which discloses determining a health indicator of a component using condition indicators (CIs) to parameterize characteristics about a component. The CI's are based on sensed data for the component.
Using wireless rather than fixed wired connections for aircraft sub-systems offers many advantages, such as lower weight, reduced complexity of wiring interfaces, simplified installation, and easier maintenance. For example, retrofitting new sub-systems into existing aircraft may become substantially easier and cheaper if no new wiring has to be designed and installed. The use of a large number of sensors to monitor aircraft's state may become more realistic if it does not lead to an increased amount of wiring. However, for wireless sensors to become a viable solution for aircraft and other applications, stringent performance requirements may need to be satisfied. Aircraft applications typically require very high reliability, i.e. data losses have to be very infrequent. The required probability of a lost data item may be several orders of magnitude lower than that typically considered by vendors of commercial wireless components. In addition, data latency requirements may be quite strict. Another issue to be addressed is the desire to limit battery replacement to very infrequent occasions when scheduled aircraft maintenance is performed, which makes energy saving important.
Accordingly, there is a need to develop solutions that address the issues noted above.


A communications system and method for self-regulating quality of service in a communications network are provided which includes measuring the amount of jitter of a first media stream received in a router, comparing the amount of jitter of the first media stream received in the router to the amount of jitter of at least one other media stream received in the router and prioritizing the timing of the transmission of a packet in the first media stream from the router based at least in part on the results of the comparing step. The communication system and method for self regulating quality of service also provide for self-regulating quality of service in a communications network by receiving at a first router a measurement of the amount of jitter associated with the media stream at a third router and prioritizing the timing of the transmission of a packet in the media stream from the first router via a second router toward the third router based at least in part on the amount of jitter associated with the media stream measured at the third router.

Current packet-based communications systems utilize packets that make up a media stream that passes through routers and/or switches as the media stream traverses from its point of origin endpoint to its final destination endpoint. Each router/switch through which the packets pass potentially introduces jitter to the packets of the media stream. Different types of media streams are prioritized as a function of their susceptibility to degradation based on jitter. The most time sensitive media streams typically are the most susceptible to degradation from jitter. For example, voice packets are typically more time sensitive than video packets, which, in turn, are typically more time sensitive than electronic mail packets. As a result, the priority of the packets will be marked accordingly. Packets are given a priority by the originating endpoint. It is the originating point that assigns the priority to the packets. There could be however a scenario where routing or switching is done in such a way that the router knows that the packet contains voice versus video versus email. The routers and switches receiving these types of packets will typically be configured to give higher priority to the voice packet over the video and the electronic mail packets while the video packet will likely be given a higher priority than the electronic mail packet. Thus, packets at a router awaiting transmission from the router will be scheduled for transmission based on the priority associated with the type of each packet awaiting transmission.

This patent is subject to a terminal disclaimer.

Network with self regulating quality of service (QoS)


A liquid distribution system having at least one liquid conduit extending from a liquid source to a liquid tap and a method for substantially retaining the temperature of a liquid in the liquid distribution system. When a tapping operation is finished, the liquid is evacuated from the liquid conduit, and a gas is brought into the liquid conduit in order to replace the liquid therein and cause the liquid to flow backwards to the liquid source. When liquid is to be tapped from the liquid tap, the gas is evacuated from the liquid conduit.

In most buildings, there is a distribution system for water, normally with separate piping for cold and hot water. Often, water is provided from a public facility. For cold water, the water is lead directly from the source, possibly after reducing the pressure, and for the hot water, there is typically a local heater or a heat exchanger, so that cold and hot water can be tapped at various locations within the building.
From the source, the conduits are often quite long, so that fresh, relatively cold water will be obtained only after tapping quite a large volume of water that has been left in the conduits when water was tapped previously, maybe several hours or even days beforehand.
A similar problem is present for hot water, which will be cooled off when there is no longer any flow of hot water in the piping system from the heating source to the respective tap. In a prior art system, as disclosed in U.S. Pat. No. 5,944,221 (Laing), there is provided a re-circulation assembly. The cooled-down (previously hot) water in the hot water line is re-circulated by a pump back to the hot water tank, where it is reheated before being delivered again. Hereby, the heat energy contained in the hot water will be fed back to the hot water tank. However, this requires a special re-circulation piping system, which is generally rather costly to install and maintain. Because of the risk of growth of micro-organisms in a hot water line, it is not permitted to re-circulate the cooled-down water through the piping intended for cold and fresh water.

Tap liquid savings in a liquid distribution system


The present invention leverages iterative transformations of search query strings along with statistics extracted from search query logs and/or web data to provide possible alternative spellings for the search query strings. This provides for spell checking that can be influenced to provide individualized suggestions for each user. By utilizing search query logs, the present invention can account for substrings not found in a lexicon but still acceptable as a search query of interest. This allows for a higher quality proposal for alternative spellings, beyond the content of the lexicon. One instance of the present invention operates at a substring level by utilizing word unigram and/or bigram statistics extracted from query logs combined with an iterative search. This provides substantially better spelling alternatives for a given query than employing only substring matching. Other instances can receive input data from sources other than a search query input.

Interaction with automated programs, systems, and services, has become a routine part of most people's lives—especially with the advent of the Internet. Web surfing or browsing for instance may even be the "new" national pastime for a certain segment of the population. In accordance with such systems, applications such as word processing have helped many become more efficient in their respective jobs or with their personal lives such as typing a letter or e-mail to a friend. Many automated features have been added to these applications such as tools for formatting documents in substantially any desired font, color, shape, or form. One tool that has been appreciated and well received by many users is a spell checking application that is either invoked by a user from the word processor to check all or portions of a respective document and/or invoked to run in the background to check spelling as users are typing. Generally, in order to perform accurate spell checking, a dictionary of "valid strings" may be employed by the spell checking application. If a spell checker encounters a string not in the dictionary, it may hypothesize that the string is a spelling error and attempt to find the "closest" string in the dictionary for the misspelled string. Most spell checkers provide a list of possible matches to the user, whereby if the match is on the list, the user can select the word having the corrected spelling from the list. Other spell checking features may perform automatic corrections—if so configured by the user.
Spell checking for word processing, however, presents only a partial view of potential areas that may be applicable to assist users when entering information into a file or document. For example, with all the potential web sites and services available, users often navigate between sites by explicitly typing in all or portions of the site name or by performing searches on words or phrases that appear in the title and the body of a web page. As many have come to find out, if the site information or the search query is entered incorrectly, the cost in time to re-navigate can become quite high. Language processors employed in search engines or other applications often process user queries and may attempt to distinguish actual user commands from incorrectly entered information. As can be appreciated, however, the type of information that may be entered for a query to a search engine may be quite different in structure or form than typically employed in a word processing application. Thus, tools that check words on a somewhat individual and isolated basis in a word processor application may have little or no utility when applied to information generated from general query data.
Browser or other search queries for information present a unique problem for spell checking applications, since the queries often consist of words that may not be found in a standard spell-checking dictionary, such as artist, product, or company names. Another problem is that a word in a query may have been entered incorrectly, but not be spelled incorrectly (for example, "and processors" instead of "amd processors"). Thus, the manner in which people enter text into a type-in line, for example, such as an input box to a search engine, is often very different than typing for word processing. Both what is entered and the types of errors people make with respect to query input are also quite different in nature. Furthermore, web data and search queries are very dynamic in nature, containing a large number of proper nouns; new products, people, institutions, locations, and events become popular every day. As such, a standard dictionary, while suitable for spell checking in the context of word processing, may not be appropriate for type-in-line and search-query spell checking.
A dictionary (i.e., lexicon) is an important component of any spell checker since the information contained therein provides the foundation to determine incorrect spellings. However, for many applications where spell checking is desired (e.g., text input provided to input boxes), a standard dictionary is not optimal for the problem. For instance, to spell check text input to the input box of a search engine, a dictionary should include strings such as "hanging chad" and "Apolo Anton Ohno" in order to check more recent events or information that may be of interest. As can be appreciated, these and a plurality of other type strings would not appear in a standard dictionary. One possible approach is to utilize substring matching techniques on a log of what users are typing into a particular location, such as a search engine or language processor. Unfortunately, a problem with this approach is that the query logs will generally also contain a large number of input errors and return substring matches that are not relevant to a user's desired search.
Additionally, the dictionary utilized for the spell checking and the context of the search are always changing. These dynamic behaviors cannot be accounted for utilizing traditional dictionary and search query processing. For example, if there is currently a popular band called Limp Bizkit, a search for "bizkit pictures" is likely to refer to this band and not a misspelling of "biscuit." If suddenly the band becomes unpopular, and there is a top-selling book on pictures of biscuits, "bizkit pictures" is then more likely a misspelling of "biscuit pictures." Likewise, given a current state of politics, "govenor anld" probably refers to "governor arnold" if he is currently a popular California governor. Thus, the context of the search query impacts the spell checking significantly.

Systems and methods for improved spell checking


A system and method for a die-to-die interconnect interface and protocol for stacked semiconductor dies. One preferred embodiment comprises an integrated circuit (IC) package comprising a first semiconductor die that includes an interface to a memory-mapped device, a second semiconductor die that does not include an interface to a memory mapped device, and a data bus coupling the first semiconductor die to the second semiconductor die (the data bus used to transfer a control word and a data word). The control word comprises a data word start address that corresponds to a location in the memory-mapped device. The data word is transferred from the second semiconductor die to the first semiconductor die and is stored by the first semiconductor die at the location in the memory-mapped device. Both semiconductor dies are mounted within the IC package.

1. Technical Field
The present subject matter relates generally to integrated circuits (ICs). More particularly, the present subject matter relates to an interconnect protocol for use in an IC package comprising at least two stacked dies.
2. Background Information
In some applications, it is desirable to include more than one semiconductor die in a single semiconductor package. For example, two dies can be stacked in a single package. As such, one die resides on top of another die. One or more issues arise regarding the interconnection protocol between stacked dies. If the dies have different "footprints" (i.e., are of different sizes), electrically interconnecting the dies may be problematic. In general, the electrical connectivity to/from a die is generally from "pins" formed or attached to the sides of the die. If the die has four sides (most dies are rectangular and thus have four sides), then all four sides of the die may be used for interconnection pins. Pins of one die can be connected to a corresponding pin of another die if the sides of the dies align vertically. For example, if two stacked dies have the same footprint, then the four sides of the dies are in vertical align and the pins on all four sides of each die can be connected with pins of the other die. However, if the dies have different footprints, then, in general, the dies can be stacked only in such a way that one or two sides (not all four) of the dies vertically align. In this case, only the one or two aligned sides can be used for die-to-die interconnection pins. As a result, fewer pins are available to interconnect stacked dies of different sizes than stacked dies of the same size.
Despite different size stacked dies being inter-connectable with fewer pins than same size stacked dies, it still may be desirable for the smaller die to have shared access to resources coupled to the larger die. These resources may have a large number of signals, well in excess of the number of interconnect pins available to the smaller die. Further, it may be desirable to share such resources without causing substantially higher power consumption of the system despite the desire for higher performance and lower latency in the face of fewer interconnecting pins.

Die-to-die interconnect interface and protocol for stacked semiconductor dies


A medical knife has a cutting edge formed along the outer edge of a blade portion, and the thickness of the lower portion located under a boundary acting as a surface including the cutting edge is made thinner than that of an upper portion located on the boundary, or a lower surface is composed of the surface including the cutting edge. The surface of the upper portion on the above surface including the cutting edge is composed of a first slanting surface constituting the cutting edge, a second reflecting surface, and a flat surface composed of a non-reflecting surface as well as the lower portion has a lower slanting surface 10 and a lower flat surface 11 each composed of a reflecting surface. With the above arrangement, the medical knife can reduce the amount of reflected light when an operation such as an ophthalmic operation is executed under a microscope.

1. Field of the Invention
The present invention relates to a medical knife which prevents a doctor from feeling the glare of illumination light used in an operation by reducing the reflection of the light as well as prevents deterioration of the cutting quality of the knife in the operation.
2. Description of the Related Art
Microscopic operations represented by an ophthalmic operation are basically executed by observing an affected area while enlarging it under a microscope. In this case, a problem arises in that it is difficult for a doctor to execute the operation because the light, which illuminates the affected area, is reflected on the surface of a medical knife and makes it difficult for the doctor to observe the affected area. To solve the above problem, there has been provided a medical knife which permits a doctor to easily observe an affected area by composing the entire surface of the knife of a non-reflecting surface. The non-reflecting surface is composed of a surface on which fine concavities and convexities are formed so that the light, which is directly reflected to the eyes of the doctor, can be suppressed by diffusely reflecting the light. In the medical knife arranged as described above, however, when the knife pierces living body tissue, friction is increased between the concavities and the convexities and the living body tissue, and further a problem also arises in that the sharpness of the knife for dissecting the living body tissue is lost because concavities and convexities are also formed on a cutting edge, and thus the cutting quality of the knife is deteriorated. When the cutting quality is deteriorated, it is difficult for an operator to form an incision wound which in agreement with his or her image, that is, it is difficult for the operator to form an incision wound having a self-closing property for permitting the incision wound to close by itself. This may adversely affect the curing of the incision wound after the operation.
Under the above circumstances, generally doctors select a medical knife depending on whether they put high priority on the cutting quality of the knife by enduring reflection of light or put high priority on reduction of reflected light by enduring deterioration of the cutting quality of the knife.
To solve the above problem, there has been proposed surgical knives whose surface, which is caught in a field of view when they are used in an operation under a microscope, is composed of a portion subjected to a reflection preventing treatment and a portion not subjected to the treatment (for example, refer to Japanese Patent Application Laid-Open (JP-A) No. 2001-57980). In the above surgical knives, the lower surface of a cutting edge formed around the periphery of the knife is entirely subjected to the reflection preventing treatment except a slanting surface constituting the cutting edge, and the upper surface of the cutting edge is composed of a slanting surface constituting the cutting edge and a roughly ground surface formed along the slanting surface constituting the cutting edge inwardly of it, and a central flat surface, which continues to the inside and the back side of the roughly ground surface and is composed of a surface subjected to the reflection preventing treatment.
In the knives according to the art disclosed in JP-A No. 2001-57980, since only the slanting surface constituting the cutting edge is composed of the reflecting surface, no light is reflected from the surfaces other than the slanting surface constituting the cutting edge. Accordingly, a workload on a doctor can be reduced because the amount of light, which is reflected when an affected area is dissected, is reduced. However, even the knives arranged as described above are not yet perfect, and it has been required to develop a medical knife which reduces an amount of reflected light and has excellent cutting quality.


This specification describes a method of visually displaying electrocardiogram data in a compressed manner on the display screen wherein rhythmic information is visible and a method of categorizing zones of the electrocardiogram data.

Electrocardiography (ECG) is the process of recording the electrical activity of the heart over a period of time, and is typically performed by applying electrodes to the patient's skin. Computerized electrocardiography produces electrocardiogram data having magnitude values of the recorded electrical activity varying over time. Electrocardiogram data can include absolute time values associated with each one of the magnitude values. Alternately, the time reference can be implicit. For instance, if the period of time elapsed between successive ones of the magnitude values is constant and known, it may not be necessary to store an absolute value of time for every magnitude value to be able to reconstruct and display the electrocardiogram.
For decades, electrocardiograms were printed directly on a paper sheet in the form of a two-dimensional graph having the magnitude plotted along the Y-axis and the time plotted along the X-axis. FIG. 1A illustrates a given period of time of an electrocardiogram presented with this method, and includes a first heartbeat 10 and the beginning of a second heartbeat 10′. A consecutive sequence of normal heartbeats are referred to as normal sinus rhythm and have reference features which are recognizable. The exact expression of these recognizable features can vary significantly from one heartbeat to another, from one person to another, and can also depend on the recording method, but are typically present in one form or another. These recognizable features include "peaks" referred to as P, Q, R, S and T peaks in the art.
In the context where normal resting heart rate typically ranges between 60 and 100 beats per minute, it will be understood from looking at FIG. 1A that electrocardiograms spanning any significant amount of time can require an imposing amount of paper.
In recent years, it has become increasingly common to store electrocardiogram data in a memory readable by a computer, and to display the electrocardiogram data on a computer display screen rather than print it on paper. Typically, only a relatively small portion of the electrocardiogram data is displayed at any given time, and a user can navigate the electrocardiogram data by providing inputs which are used as commands to change the displayed portion of the electrocardiogram data.
There was a need for displaying larger periods of time of electrocardiogram data at once on the display screen, e.g. minutes of electrocardiogram data. This need was met to a certain extent by a new method of displaying electrocardiogram data which has been published in recent years. This new method involved compressing the displayed data by using a color scale to plot magnitude rather than by using the Y-axis of the 2D graph. Accordingly, the display of the electrocardiogram data could be compressed to a single axis. Moreover, algorithms have been developed to recognize individual heartbeats automatically in the electrocardiogram data. Accordingly, single-axis graphs could be displayed for corresponding heartbeats. By orienting the single-axis graphs transversally and displaying them side-by-side, an impressive amount of heartbeats could be displayed on the display screen all at once.
While former technologies were satisfactory to a certain degree, there remained room for improvement in terms of displaying, navigating, annotating, and otherwise handling the electrocardiogram data using a computer.

Computer-implemented method of handling electrocardiogram data


A system and related method of operation for migrating the memory of a virtual machine from one NUMA node to another. Once the VM is migrated to a new node, migration of memory pages is performed while giving priority to the most utilized pages, so that access to these pages becomes local as soon as possible. Various heuristics are described to enable different implementations for different situations or scenarios.

Modern computer systems provide support for virtual memory. With this support, an application sees a contiguous address space that is not tied to underlying physical memory in a computer system. In particular, an application accesses a virtual address (VA) that is translated to a physical address (PA) which is used to access the physical memory. The translation proceeds in one of two ways, depending on whether the system architecture supports "segmentation" or not. If the system architecture supports segmentation, then a VA is added to a value in a "segment base register," to produce what is called a linear address (LA). The value of the VA is compared with a value in another register called a "segment limit register." If the value of the VA is greater than the value in the segment limit register, then a "general protection fault" is generated, and execution jumps to a fault handler. Otherwise, the LA is translated to the PA using hardware called a Memory Management Unit (MMU). If the system architecture does not support segmentation, then the LA is the same as the VA, and the VA is used by the MMU to translate to the PA.
FIG. 1B is a block diagram that helps illustrate how the MMU performs the translation process. The following uses a 64 bit architecture and a 4 KB page size for sake of illustration. However, similar steps may be utilized for a 32 bit architecture, and different page sizes may be utilized as well. A more detailed discussion of translation processes can be found in "Intel 64 and IA-32 Architecture Application Note: TLBs, Paging-Structure Caches, and Their Invalidation" available from Intel of Santa Clara, Calif. Although this specification describes translation processes for Intel processors, the details are similar for other popular x86 based processors, and embodiments of the invention may be implemented on any x86 or non-x86 based processor that supports virtual memory.
With respect to the x86 architecture, the MMU consists of the following parts: a control register CR3, a translation look aside buffer (TLB), and translation circuitry. The paging structures: (a) are 4 KB in size; (b) reside in main memory, generally separate from the MMU; and (c) are designated L4, L3, L2 and L1 page tables (they are sometimes also called level 4, 3, 2 and 1 page tables), respectively. In some implementations, the page table structures contain 512 8-byte entries comprising information required to perform the translation. The following information is relevant here. First is a CR3 register which contains the physical address of some L4 page table. Second is a physical page number of the next lower level page table. For example, referring to FIG. 1B, a page table entry in the L4 page table contains the physical page number of the L3 page table. Third is a present bit (P-bit). The P-bit indicates whether a page table entry is "present" or "not present." If the P-bit is 0 in any page table entry, translation stops and a page fault is generated. Fourth is an access bit (A-bit). When a linear address is used to access memory, the processor sets the A-bit to 1 in all page table levels used to translate the linear address.
As is well known, translation of a linear address LA (generally referred to as "walking the page table" or a "page table walk") proceeds as follows. A linear page number (LPN) is formed from the most significant 52 bits of the LA. The least significant 12 bits of the LA constitute what is called an "offset." Using the CR3 register, the physical address of the L4 page table is determined. An entry in the L4 page table is selected by bits 47:39 of the LA; the selected entry contains a physical address of the L3 page table. An entry in the L3 page table is selected by bits 38:30 of the LA; the selected entry contains a physical address of the L2 page table. An entry in the L2 page table is selected by bits 29:21; the selected entry contains a physical address of the L1 page table. An entry in the L1 page table is selected by bits 20:12; the selected entry contains a PPN corresponding to the LPN. Finally, the PA comprises the PPN as its 52 most significant bits, and the offset as its 12 least significant bits.
As is well known, an MMU is typically equipped with one or more TLBs, where a TLB is a cache of recent LPN to PPN translations. To translate an LA, the MMU computes the LPN, and then looks in the TLB for a translation of the LPN. If the translation is present in the TLB (referred to as a "TLB hit"), the PPN is immediately available from the TLB. On the other hand, if the translation is not present in the TLB (referred to as a "TLB miss"), a page table walk is done, and the translation is stored in the TLB, potentially evicting another entry from the TLB.
As is well known, virtualizing an MMU so that multiple virtual machines can run on a single hardware system typically entails another level of translation. The first translation is provided by a guest operating system (guest OS) running in a virtual machine. The guest OS translates a guest LPN (GLPN) into a corresponding guest PPN (GPPN) in the conventional manner. The second translation is provided by virtualization software, for example, a virtual machine monitor (VMM). In particular, the VMM maintains a GPPN to "host PPN" mapping in its internal translation table where the host PPN is used to address physical memory of the hardware system. In the rest of this specification, this translation table will be referred to as T.
One of two methods is typically used for virtualizing an MMU, namely a shadowing of guest paging structures method (shadowing method), and a hardware assist method. FIG. 1C illustrates the shadowing method. As shown in FIG. 1C, in the shadowing method for virtualizing an MMU, virtualization software, for example, the VMM, maintains shadow page tables, one shadow page table for each guest page table. While the guest page tables, maintained by the guest operating system, contain guest LPN to guest PPN mappings, the shadow page tables contain guest LPN to host PPN mappings. To insert a translation for a guest LPN gl into a shadow page table, the VMM walks the guest page table to determine the guest PPN gp. Then, it translates the guest PPN gp to a host PPN hp using its translation table T.
In operation, for a particular guest, the VMM changes the CR3 register so that it points to the shadow page table for that guest. By doing this, the hardware MMU uses the shadow page table for address translation. Further, the VMM saves the value of the CR3 register for that guest. To ensure correct guest operation, the VMM traps guest accesses to the CR3 register. Thus, whenever the guest reads the CR3 register, the VMM provides it with the saved value of the CR3 register. Further, whenever the guest writes to the CR3 register, the VMM saves the new value of the CR3 register for the guest, and writes to the physical CR3 register the physical address of the new shadow page table.
When a new shadow page table is loaded by the VMM (by writing a new value to the CR3 register in response to a guest writing to the CR3 register), it is not necessary that the shadow page table contains valid entries corresponding to all valid entries in the guest page table. For a guest page table entry, if the PPN in its corresponding shadow page table entry is invalid, the VMM sets the P-bit of that shadow page table entry to 0. Then, when this shadow page table entry is used to translate an address, a page fault will be generated. Such a page fault is called a "hidden page fault" because it would not have been generated if the guest operating system were executing in a non-virtualized environment. In response, a page fault handler of the VMM determines the faulting guest linear page number gl (using a faulting virtual address put into a predetermined hardware register), and translates it to a guest physical page number gp by walking the guest page tables. Then, the guest physical page number gp is translated to a host physical page number hp using translation table T. Finally, the host physical page number hp is inserted into the concerned shadow page table entry, the P-bit of the entry is set to 1, and the faulting guest instruction is restarted.
The VMM needs to ensure that shadow page tables are synchronized with guest page tables. For example, if the guest OS changes a translation of guest linear address gl from gp to gp′, the VMM needs to update the corresponding entry in the shadow page table as well.
The hardware assist method for virtualizing an MMU makes use of a virtualization assist provided in certain modern CPUs. As is well known, apart from conventional page tables provided by earlier CPUs, certain modem CPUs provide a second layer of page tables. The design of the second layer of page tables, called "hypervisor page tables," is similar to that of the guest page tables as described above. The guest page tables, maintained by the guest operating system, continue to contain guest linear to guest physical translations. However, the hypervisor puts the guest physical to host physical translations in the second layer of page tables (description of hypervisor is provided in the Appendix below). Apart from the CR3 register provided by earlier CPUs, the virtualization assist of the modern CPUs provides an hCR3 register which contains the host physical address of the L4 page of the hypervisor page tables.
With the hardware assist, on a TLB miss, the MMU walks through both sets of page tables to translate a guest linear address to a host physical address. This is done by using translations from the guest page tables and the hypervisor page tables: the guest linear address is translated to a guest physical address using the guest page tables, and the guest physical address is translated to a host physical address using the hypervisor page tables. Using the hardware assist obviates the construction and maintenance of shadow page tables.
FIG. 1D illustrates an example of virtual machine memory virtualization. As shown in FIG. 1D, each of the virtual machines VM1 and VM2 runs two processes: process 1 and process 2. The processes execute using virtual memory 122, which is mapped to physical memory 124, and physical memory 124 is mapped, in turn, to machine memory 126. The guest OS in each virtual machine maps virtual memory (i.e., virtual page numbers (VPN)) to guest physical memory (i.e., guest physical page number (GPPN))—this is shown by inset 112. In addition, virtualization software, for example, a VMM, maintains a mapping of guest physical page numbers (GPPN) to actual machine or host page numbers (MPN)—this is shown by inset 114. The VMM also maintains a mapping of VPN to MPN in shadow page tables to map virtual memory directly to machine memory (i.e., VPN to MPN) to avoid two levels of translation on every access—this is shown by inset 116. When the guest OS changes the virtual memory to physical memory mapping, the VMM updates (or invalidates) the shadow page tables to enable a consistent lookup (or generation of a hidden page fault).
Various advancements in microprocessor fabrication have led to modern CPUs that are significantly faster than main memory. Consequently, CPUs running demanding applications may become constrained by memory bandwidth. This issue may be exacerbated in symmetrical multiprocessing (SMP) systems wherein applications running on several CPUs access main memory together, thereby competing for access to a shared memory bus. To overcome this issue, a Non-Uniform Memory Access (NUMA) architecture has been developed.
In NUMA machines, processors are divided into NUMA nodes where each node has a separate memory, and a bus connects the processors of a node to the memory. An example is illustrated in FIG. 1E where each NUMA node, NUMA 101 and NUMA 102, has several processors coupled to local memory. As further shown in FIG. 1E, the different NUMA nodes are connected by an interconnect 103. In NUMA machines, each processor can access any memory, including memory on different NUMA nodes; however, access to memory on the same NUMA node (referred to as "local memory") is faster than access to memory on a different NUMA node (referred to as "remote memory"). Due to this, software designers attempt to ensure that data most often accessed by a program running on a CPU (called a "working set" of the program) usually reside on the same NUMA node as the CPU.
From time to time, it may become necessary or desirable to migrate a virtual machine from one NUMA node to another. Consider, for example, a situation illustrated in FIG. 1E where NUMA node 101 includes processor A having two virtual machines, VMA1 and VMA2, running demanding applications, thereby heavily loading processor A. On the other hand, assume that one or both of processors A′ and B′ on NUMA node 102 are lightly loaded. To alleviate the load on processor A, it may be desirable or necessary to migrate one of VMA1 and VMA2 to NUMA node 102. Such load balancing may be done automatically by continuously monitoring the load on all of the nodes.

System and method for improving memory locality of virtual machines


A plurality of memory cell devices is formed by using an intermediate dual polysilicon-nitride control electrode stack overlying nanoclusters. The stack includes a first-formed polysilicon-nitride layer and a second-formed polysilicon-containing layer. The second-formed polysilicon-containing layer is removed from areas containing the plurality of memory cells. In one form the second-formed polysilicon-containing layer also contains a nitride portion which is also removed, thereby leaving the first-formed polysilicon-nitride layer for the memory cell devices. In another form the second-formed ploysilicon-containing layer does not contain nitride and a nitride portion of the first-formed polysilicon-nitride layer is also removed. In the latter form a subsequent nitride layer is formed over the remaining polysilicon layer. In both forms a top portion of the device is protected from oxidation, thereby preserving size and quality of underlying nanoclusters. Gate electrodes of devices peripheral to the memory cell devices also use the second-formed polysilicon-containing layer.

1. Field
The present invention relates to semiconductor devices, and, more particularly, to such devices having nanoclusters.
2. Description of the Related Art
Some devices such as memories (e.g. non volatile memories) utilize discrete charge storage elements called nanoclusters (e.g. of silicon, aluminum, gold, or germanium) for storing charge in a charge storage location of a transistor. In some examples, the nanoclusters are located between two dielectric layers, a bottom dielectric and a control dielectric. Examples of such transistors include thin film storage transistors. A memory typically includes an array of such transistors. Examples of nanocluster types includes doped and undoped semiconductor nanoclusters such as silicon nanocrystals, germanium nanocrystals and their alloys. Other examples of nanocluster types include various conductive structures such as metal nanoclusters (e.g., gold nanoclusters and aluminum nanoclusters), and metal alloy nanoclusters. In some examples, nanoclusters are from 10–100 Angstroms in size.
Some memories that have charge storage transistors with nanoclusters are implemented on integrated circuits that also include high voltage transistors in the circuitry used for charging and discharging the charge storage locations of the charge storage transistors. Charging or discharging the charge storage locations is used to store one or more bits of information, and may be referred to as programming or erasing. These high voltage transistors typically include a relatively thick gate oxide. This gate oxide may be formed under severe oxidizing conditions. This oxidizing ambient may penetrate the control dielectric of the charge storage transistors thereby undesirably oxidizing the nanocrystals and undesirably increasing the bottom dielectric thickness. Accordingly, an improved method for making a device with nanoclusters is desirable.
The use of the same reference symbols in different drawings indicates similar or identical items. Skilled artisans appreciate that elements in the figures are illustrated for simplicity and clarity and have not necessarily been drawn to scale. For example, the dimensions of some of the elements in the figures may be exaggerated relative to other elements to help improve the understanding of the embodiments of the present invention.

Method of forming a nanocluster charge storage device


A circular filtering system that prevents the problem of inter-symbol interference. The circular filtering system utilizes a buffer memory to store samples of a given symbol and provide only these samples to a linear filter such that the output of the filter, for any given symbol is formed by filtering only samples of that input symbol. Each symbol being filtered independent of other symbols hence eliminating inter-symbol interference caused by filtering. Where symbols are tolerant to a fixed phase shift for each symbol, the circular filtering system can be simplified by reducing the size of the buffer and introducing a multiplexer.

1. Field of the Invention
The invention relates generally to digital signal processing. More specifically, the invention relates to processing digital signals that are made up of concatenated sub-sequences.
2. Related Art
In systems that perform digital signal processing, there are classes of these systems which perform signal composition and decomposition. In the case of signal composition, digital signals are being synthesized from a set of data or from analog signals. In the case of signal decomposition, digital signals are often received by a system so that they can be interpreted, analyzed, or merely detected. For instance, a transmitter may digitally model a set of data into digital signals. Then, a receiver, would estimate the contents of the original set of data by analyzing the digital signals that are received.
Certain signals that are at or near a zero frequency (i.e. DC) are called "baseband" signals. FIG. 1 shows one such baseband signal 100 centered around a frequency of zero. Baseband signal 100 has an upper limit in frequency fsignal-max at which the signal level drops to zero in amplitude and has no level higher than that at higher frequencies. A discrete-time signal has an infinite number of "images" centered around higher frequencies which are multiples of a sampling frequency fsampling. The baseband signal 100 has images 110, 120 and 130 that occur at discrete multiples of fsampling which are 1*fsampling, 2*fsampling, and 3*fsampling, etc., correspondingly, since the baseband spectrum is centered around a frequency of zero. There are other images at higher multiples, that are not shown.
The digital signal (composed of baseband signal 100 and its images 110, 120 and 130) shown in FIG. 1 is a signal that does not exhibit a condition known as "aliasing". Aliasing occurs when a baseband signal is processed (for instance, up-sampled) such that images overlap one another in their frequency spectrums. If the sampling frequency at which a baseband signal were up-sampled were low enough, then images might overlap partially or even entirely. Referring to FIG. 1, if fsampling were less than 2*fsignal-max, and assuming the baseband signal 100 is symmetric about zero frequency (amplitude axis), then the trailing edge of image 110 would overlap with the baseband signal 100. The signal would then be aliased, and thus corrupted such that the original signal could not easily, if at all, be recovered (since the original content of the baseband is a priori unknown in most cases). One way of preventing aliasing is to choose a high enough sampling frequency which can make aliasing impossible. Signals that exhibit no aliasing, such as that shown in FIG. 1, have "guard bands" between the leading edge of one signal image and trailing edge of the next signal image. For instance, between baseband signal 100 and image 110, there is a guard band between them of:guard band=fsampling−2*fsignal-max.
Signals and systems that have a large (relative to the total spectrum size of the baseband signal) guard bands are often used in applications where the prevention of aliasing is crucial. FIG. 2(a) is an example of a digital signal with a larger guard band. Baseband signal 200 is shown having images at multiples of a sampling frequency f′sampling that is much larger than fsampling shown in FIG. 1. Such signals with large frequency guard bands are used in systems that may involve many other signal processing steps. The size of the guard band can often simplify these other steps. For instance, systems which convert digital (or discrete-time) signals into analog signals (continuous time) would have to deal with images of the signal. If there were large guard bands, it would be easier to deal with such images (since there is more room between images).
In the time-domain, such signals shown in the frequency domain as waveforms may actually be composed of or encode "symbols". When analyzed as symbols in the time-domain, "inter-symbol interference"(ISI), aliasing becomes a major issue. ISI is interference between two or more symbols contained within a given signal. For example, if a digital signal x(n) is composed of a series of N-length concatenated sub-sequences (i.e. symbols) xi(n), thenx(n)={x0(0), x0(1), . . . , x0(N−1), x1(0), x1(1), . . . , x1(N−1), . . . }.In this example, ISI involves interference between a symbol x0 and the next symbol x1. In this example, ISI is caused by the use of linear filters that operate to output one result for every cycle. Such linear filters have a designated number of taps, M, each of which serve to multiply a provided input by a filter coefficient and sum all such results over the M taps. Linear filters achieve this by well-known sample-and-hold or delay mechanisms which have the effect of pipelining inputs so that they are eventually multiplied by every one of the M taps.
FIG. 2(b) depicts a conventional M-tap linear filter. Initially, the M-tap filter 240 pictured takes M−1 samples of a previous input symbol and the first sample of the current symbol to produce the first sample for the output of the current symbol. Likewise, the second output sample is produced by the M-tap filter 240 filtering the last M−2 samples of the previous symbol and the first two samples of the current symbol. Thus, each symbol contributes samples to the processing of the next symbol, and thus interferes with the next symbol. Once M samples are processed, then the filter 240 takes the current sample of the same symbol and the M−1 samples of the same symbol in order to generate succeeding output samples.
Because of the sequential nature of these filters, a given output is often a filtered version of samples of one symbol and the next. For instance, consider the case of xi above, where N=10. Each symbol x0, x1 etc. would be composed of 10 samples. If the number of taps M in the filter is 6 (often, the number of taps is less than or equal to the number of samples that need processing, depending on the application), then when the first sample of symbol x1(0) is input to the filter, the last 5 samples x0(5) . . . x0(9) of the previous symbol x0 will be residual inputs remaining in the filter therewith. The first output of the linear filter for the first sample x1(0) of symbol x1 is thus corrupted and exhibits ISI (between x(1) and the sequentially preceding symbol x(0)).
This effect is further illustrated in FIG. 3(a). The input is comprised of three concatenated symbols, indexed as i−1, i, and i+1. The next two rows depict the filtered versions of symbols i−1 and i, respectively. The spreading out at the ends of each symbol into the area of adjacent symbols is caused by filtering. The final row shows the area of overlap between the pair of symbols i−1 and i, which represents where ISI occurs. FIG. 3(b) shows a simulation of the example depicted in FIG. 3(a). The input 310 consists of three concatenated symbols each of which is composed of two cycles of a sinusoid, but of differing phase. The filter (not shown) applied upon the input 310 is a low-pass filter with a cutoff frequency above the frequency of the sinusoid. Ordinarily, the waveforms should be unchanged in such a case. However, as depicted in FIG. 3(b) the output 320 which results from such filtering of the input 310, the output 320 shows distortion at the symbol boundaries. This is a result of ISI.
The conventional techniques for generating and detecting signals with large frequency guard bands often have the effect of also causing ISI, since these techniques were intended for use in signals that do not have sub-sequences of symbols contained within them. For instance, one technique for generating the signal 20 shown in FIG. 2(a) involves up-sampling the baseband signal 200 by two using the well-known zero insertion process. After up-sampling, the resulting up-sampled signal is then filtered using a digital lowpass filter 400 such as that shown in FIG. 4. In the frequency domain (shown), the sharp response (i.e. short transition region) of the low-pass filter 400 implies it will exhibit a very long impulse response in the time domain. In digital systems, this long impulse response will cause ISI (because neighboring symbols are corrupted by the length of the response).
Other techniques have sought to convert the digital signal into an analog signal and then perform filtering using an analog filter. After analog filtering the filtered signals may be converted back, at say a receiver, to digital. However, analog low-pass filters have a similarly long impulse response and thus, subject the signal to ISI which is exhibited when the signal reverts to digital form.
For some systems, the digital synthesis/analysis is performed by the well-known Discrete Fourier Transform (DFT) and/or an Inverse Discrete Fourier Transform (IDFT) calculation. The DFT is often implemented by a special transform, which is efficient to implement in digital systems, known as the Fast Fourier Transform (FFT). In these cases, the conventional length N of the FFT is the channel-width (frequency spectrum width) divided by the sub-carrier spacing. For example, one emerging standard in wireless multi-channel communications systems, IEEE (The Institute of Electrical and Electronic Engineers) 802.11a, the channel width is 20 MHz (MegaHertz) and the subcarrier spacing is 312.5 KHz (KiloHertz) which results in an FFT length of N=(20*106/312.50*103)=64. FIG. 5 shows a spectrum of a typical multi-channel communications medium with 20 MHz channels. The IEEE 802.11a standard specifies that each 20 MHz channel has 52 sub-carriers covering 16.5 MHz of the 20 MHz, leaving 3.5 MHz to be used for preventing interference between channels. The FFT length is chosen conventionally to be 64-point for architecture and implementation efficiency, even though only 52 would be needed. For the synthesis (transmitter) case, the signal that results from a 64-point FFT is the same as that shown in FIG. 1. For the analysis (receiver) case, the input to the FFT has the same spectrum as shown in FIG. 1. Both cases require difficult filtering to increase their guard bands and prevent time-domain exhibited ISI.
FIG. 6 illustrates the filtering of an input symbol stream using conventional linear filtering. An input symbol stream is decomposed into a number of groups (symbols) of samples of a fixed length N, where N is the number of samples per symbol. To decode the first sample of a current symbol 620, which is indexed by variable "i", the first sample of the current symbol 620 is passed to an M-tap linear filter 640. The M-tap filter 640 also holds in its taps the last M−1 samples of the previous symbol. The resulting decoded output sample 660 of the decoded symbol 650 (labeled i′) will be the filtered result from part of the current symbol 620 (i) and a portion of the previous symbol 610 (i−1). Since the output of the current symbol's samples may be based in part on the input from a prior symbol, the intersymbol interference will contribute to an inaccurate result such that the output 660 may not be as predicted by the filter response. While additional techniques have been used to remove such interference after filtering, few fail to address the interference itself.
Thus, there is a need for methods that process digital signals that have large frequency guard bands but do so without much ISI.

Efficient methods for filtering to avoid inter-symbol interference and processing digital signals having large frequency guard bands


An improved baler and a method of using the improved baler to produce high density bales. The baler comprises a pick-up assembly configured to pick-up crop material; a compression assembly comprising at least one pair of opposing compression rollers configured to generate a pressure on the crop material sufficient to crush the nodes as it passes between the compression rollers; and a bale chamber, wherein the crushed crop forms a bale. The compression rollers define an adjustable gap therebetween. By adjusting the gap, a pressure sufficient to crush the nodes of the crop may be exerted on the crop material as it passes between the compression rollers may be controlled, resulting in bales of higher density than conventional bales.

For many years agricultural balers have been used to consolidate and package crop material so as to facilitate the storage and handling of the crop material for later use. Usually, a mower-conditioner cuts and conditions the crop material for windrow drying in the sun. When the cut crop material is properly dried, a baler, for example a round baler, travels along the windrows to pick up the crop material and form it into cylindrically-shaped round bales.
More specifically, a pick-up assembly of the baler gathers the cut and windrowed crop material from the ground then conveys the cut crop material into a bale-forming chamber (or bale chamber) within the baler. A drive mechanism operates to activate the pick-up assembly, augers, and a rotor of the feed mechanism. A conventional bale chamber may include of a pair of opposing sidewalls with a series of belts, chains, and/or rolls that rotate and compress the crop material into a cylindrical shaped bale. Square balers operate on similar principles but form "slices" which when stacked together form a rectangular or "square" bale.
Certain dry crop materials, such as miscanthus, switchgrass, small grain straws, or a mix of corn cobs and stover have a relatively low bulk density in traditional bales, either round or square. Densities can range from 7 to 14 pound/square feet. To be able to take advantage of trucking capacities, the U.S. Department of Energy has targeted 15.5 pound/square feet as a goal to improve transportation and logistic costs. Moreover, crop materials with stems having nodes and/or internodes also have a relatively low bulk density in round or square traditional bales.
Current round balers today rely on belt tension to regulate the bale density. However, as belt tension increases, the durability of the belt and rolls decreases. Thus, further increasing belt tension does not serve as a good option to increase bale density. Relatedly, it is also desirable to improve current square balers today to form high density bales.
High density bales are desirable because they would allow people to take advantage of transportation capacities and will reduce the numbers of the bales to work with. The present invention is directed to these and other important ends.

Balers and methods for forming high density bales


A system for providing secondary data storage and recovery services for one or more networked host nodes has a server application for facilitating data backup and recovery services; at least one client application for facilitating host node configuration to receive services; a secondary data storage medium; and at least one mechanism for passive acquisition of data from the one or more host nodes for storage into the secondary data storage medium by the server application. In a preferred embodiment secondary storage is streamlined through continuous data backup and enhanced by elimination of redundant write data.

In the field of data storage and recovery, storage area networks (SANs) are being used more and more frequently to store data because of their high speed, reliability, and relative fault tolerance characteristics.
A SAN is a high-speed network that separates storage traffic from other types of network traffic. Prominent standards used in data storage in conventional art are Small Computer Systems Interface (SCSI), Fibre Channel (FC), SAS and ATA/SATA. The redundant array of independent disks (RAID) standard is used for creating fault tolerant data storage.
State-of-the-art storage systems comprise a SAN network of storage devices and host nodes, in many instances connected together by a Fibre Channel (FC) switch. There are other variations of SAN architecture that use a different transport protocol such as high speed Ethernet in place of the FC protocol for the storage network.
In addition to the above, late-developing enhancements such as iSCSI and Fibre Channel over Internet protocol (FCIP) have enabled data storage networks to be managed as two or more networked SAN storage islands connected into a larger network by an IP tunnel. In these cases, which are not yet widely practiced, TCP/IP is leveraged as well as encapsulation (frame packaging) methods to enable data storage devices to communicate with each other efficiently over the Internet in a dedicated manner.
In typical application data generated by nodes on a host network is written to a primary storage on a SAN. Data from the primary storage system is then typically archived or backed up to a tape media in a batch-mode fashion at the end of a work period. Typically, a larger number of data-generating machines in a host network, like PCs and servers, back up the data to a smaller number of mass storage devices like a tape library. For many applications leveraging an off-site storage solution, data written to a primary storage system is transferred to one or more tape drive systems as described above for archiving to magnetic tape media, which can be securely stored off-location on behalf of an enterprise.
A typical problem with the backup operation (writing the data to tape) is that data generated from some machines during any work period can be of a very large volume and can take considerably longer to back-up than data from other machines. The backup window for a specific host can range anywhere from 30 minutes to 48 hours and above depending on the volume of data changes generated.
Another problem is that the backup data is sent from the host nodes to the tape drive over the LAN. Rendering the data from RAID to tape is typically done in a manually orchestrated batch mode operation performed by an administrator with the help of backup software. Under these conditions the operating host data network (LAN) must share bandwidth with the components involved in securing the backup data to tape media.
Yet another limitation with prior art systems is that if a data recovery operation is required, wherein the data desired is already archived to tape media, the recovery process is comparatively much slower than, for example, recovery of near-term data from a hard-drive disk.
There are still more limitations are apparent with practices of prior art storage and backup systems. For example, with prior art data backup software data movement is from each of the hosts. Moreover, most prior art backup systems perform backup operations at the file level in a non continuous fashion, which can cause additional disk seeks to find out which files have actually changed by the scheduled backup time.
What is needed in the art is a method and apparatus for archiving data to a backup data-storage sub-system that solves the above problems.

Secondary data storage and recovery system


A common-mode choke coil has first and second coil parts, a first inductance part, and a second inductance part. The first coil part and the second coil part are magnetically coupled to each other. The first inductance part is electrically connected in series to the first coil part and is not substantially magnetically coupled to the first coil part. The second inductance part is electrically connected in series to the second coil part and is not substantially magnetically coupled to the second coil part. The first inductance part and the second inductance part are not substantially magnetically coupled to each other.

1. Field of the Invention
The present invention relates to a common-mode choke coil.
2. Related Background Art
A known common-mode choke coil of this kind includes first and second coil parts magnetically coupled to each other (e.g., cf. Japanese Patent Application Laid-Open No. 8-203737). In the common-mode choke coil described in Japanese Patent Application Laid-Open No. 8-203737, a coil conductor constituting the first coil part and a coil conductor constituting the second coil part are located between a pair of magnetic substrates.


An airfoil for a stator vane having an uncoated profile substantially in accordance with Cartesian coordinate values of X, Y and Z set forth in Table I is provided. The profile is carried only to three decimal places wherein Z is a distance from a platform on which the airfoil is mounted and X and Y are coordinates defining the profile at each distance Z from the platform.

The present invention relates generally to stator vanes for gas turbines and, more particularly, to a novel and improved profile for a ninth stage compressor stator vane.
In the design, fabrication and use of turbine engines, there has been an increasing tendency toward operating with higher temperatures and higher operating pressures to optimize turbine performance. Also, as existing turbine airfoils and stator vanes reach the end of their life cycle, it is desirable to replace the airfoils, while simultaneously enhancing performance of the gas turbine through redesign of the airfoils to accommodate the increased operating temperatures and pressures.
Airfoil profiles for gas turbines have been proposed to provide improved performance, lower operating temperatures, increased creep margin and extended life in relation to conventional airfoils. See, for example, U.S. Pat. No. 5,980,209 describing an enhanced turbine blade airfoil profile. Advanced materials and new steam cooling systems now permit gas turbines to operate at, and accommodate, much higher operating temperatures, mechanical loading, and pressures than is capable in at least some known turbine engines. As a result, many system requirements must be met for each stage of each compressor used with the turbine engines in order to meet design goals including overall improved efficiency and airfoil loading. Particularly, the airfoils of the stator vanes positioned within the compressors must meet the thermal and mechanical operating requirements for each particular stage.


Systems and methods for measuring image recognition accuracy are provided. One embodiment of a method includes locating an identifier in an environment, where the identifier is configured according to a predetermined format and where the identifier identifies an actual characteristic of an object. Some embodiments of the method include locating the object in the environment, determining a perceived characteristic of the object, and determining, from the identifier, the actual characteristic of the object. Similarly, some embodiments include comparing the actual characteristic of the object with the perceived characteristic of the object, determining whether the actual characteristic of the object substantially matches the perceived characteristic of the object and, in response to determining that the actual characteristic of the object does not substantially match the perceived characteristic of the object, determining a change to make for improving image recognition.

Robots are being utilized for an increasing number of tasks, such as cleaning an environment, transporting objects from a source area of an environment to a destination area of the environment, etc. Regardless of the particular use of the robot, the robot may utilize image recognition to determine various features of an environment. However, it is often difficult to determine the quality of the image recognition and there are often limited resources for improving the image recognition while the robot is in use.

Toyota Motor Engineering & Manufacturing North America, Inc.


This invention provides compounds of Formula (I), having the structure

Protein kinases are a family of enzymes that catalyze the transfer of the terminal phosphate of ATP the hydroxyl group of specific tyrosine, serine, threonine, or histidine residues in protein. It is known that such phosphorylation plays a fundamental role in essentially all molecular aspects of cell life including metabolism, cell proliferation, cell differentiation, cell migration, and cell survival, and that protein kinases constitute major pharmacological targets [Schlessinger and Ullrich, Neuron, 9, 383 (1992); Cohen, P. Nat. Rev. Drug Discov. 1, 309-315 (2002); Scapin G., Drug Discovery Today 7(11): 601-611 (2002)].
Of the many different cellular functions in which the activity of protein kinases is known to be required, some processes represent attractive targets for therapeutic intervention for certain disease states. For example, specific protein kinases have been implicated as targets in cancer [Traxler, P. M., Exp. Opin. Ther. Patents, 8, 1599 (1998); Bridges, A. J., Emerging Drugs, 3, 279 (1998)], restenosis [Mattsson, E., Trends Cardiovas. Med. 5, 200 (1995); Shaw, Trends Pharmacol. Sci. 16, 401 (1995)], atherosclerosis [Raines, E. W., Bioessays, 18, 271 (1996)], blood vessel proliferative disorders such as angiogenesis [Shawver, L. K., Drug Discovery Today, 2, 50 (1997); Jackson et al J. Pharm. Exp. Ther. 284, 687 (1998); Folkman, J., Nature Medicine, 1, 27 (1995)], chronic obstructive pulmonary disease, bone disease such as osteoporosis [Boyce, J. Clin. Invest., 90, 1622 (1992), Tanaka et al, Nature, 383, 528 (1996)], psoriasis [(Dvir, et al, J. Cell Biol. 113, 857 (1991)], inflammatory disorders such as arthritis [(Badger, J. Pharm. Exp. Ther. 279, 1453 (1996)], central nervous system disorders such as Alzheimer's [(Mandelkow, E. M., et al, FEBS Lett, 314, 315 (1992); Sengupta, A. et al, Mol. Cell. Biochem. 167, 99 (1997)], pain sesation [Yashpal, K., J. Neurosci. 15, 3263-72 (1995)], autoimmune diseases and transplant rejection [Bolen and Brugge, Ann. Rev. Immunol. 15, 371 (1997)], thrombosis [Salari, FEBS, 263, 104 (1990)], metabolic disorders such as diabetes [Borthwick, A. C. et al, Biochem. Biopys. Res. Commun. 210, 738 (1995)], and infectious diseases (Lum, R. T. PCT Int. Appl., WO 9805335A1 980212), and viral infections [Littler, E. Nature, 160, 358 (1992)].
A partial, non-limiting, list of such kinases includes CDK1, CDK2, CDK3, CDK4, CDK5, CDK6, CDK7, CDK8, CDK9, CDK11, PDK1, PDK2, cRaf1, c-src, abl, Araf, ATK, bcr-abl, Blk, Braf, Brk, Btk, cfms, c-fms, c-kit, c-met, CSF1R, CSK, EGFR, ErbB2, ErbB3, ErbB4, ERK, ERK1, ERK2, Fak, fes, FGFR1, FGFR2, EGFR3, EGFR4, EGFR5, Fgr, FLK-4, Fps, Frk, Fyn, GSK, gsk3a, gsk3b, Hck, IGF-1R, IKK1, IKK2, IKK3, INS-R, integrin-linked kinase, Jak, JAK1, JAK2, JAK3, JNK, Lck, Lyn, MEK, MEK1, MEK2, p38, PDGFR, PIK, PKB1, PKB2, PKB3, PKC, PLK1, Polo-like kinase, PYK2, tie1, tie2, TrkA, TrkB, TrkC, UL13, UL97, VEGF-R1, VEGF-R2, Yes, AKT, and Zap70.
As an example, in cancer cells, IGFR-1 plays a critical role because it contributes to the promotion of tumor growth by inhibition of the apoptosis, transformation, metastasis and induction of angiogenesis through the vascular endothelial growth factor.
In addition to tyrosine kinases, there are serine/threonine protein kinases, that phosphorylate serine and/or threonine residues on proteins. Among them, cyclin-dependent kinases (CDKs) play a key role in regulating the cell cycle machinery. These complexes consist of two components: a catalytic subunit (the kinase) and a regulatory subunit (the cyclin). To date, 13 CDKs have been identified along with 25 cyclin-box-containing proteins [Knockaert, M.; Greengard, P.; Meijer L. Trends in Pharmacological Sciences 23(9): 417-425.] Each kinase associates with a specific regulatory partner and together make up the active catalytic moiety. Each transition of the cell cycle is regulated by a particular CDK complex: G1/S by CDK2/cyclin E, CDK4/cyclin D1 and CDK6/cyclin D2; S/G2 by CDK2/cyclin A and CDK1/cyclin A; G2/M by CDK1/cyclin B, the coordinated activity of these kinases guides the individual cells through the replication process and ensures the vitality of each subsequent generation [Science, 274, 1643-1677 (1996); Ann. Rev. Cell. Dev. Biol., 13, 261-291 (1997); Fischer, P. M. Current Opinion in Drug Discovery and Development, 4(5), 623-634 (2001); Draetta, Trends Biochem. Sci. 15:378-382 (1990); Sherr, Cell 73:1059-1065 (1993)].
An increasing body of evidence has shown a link between tumor development and CDK related malfunctions. Over-expression of the cyclin regulatory proteins and subsequent kinase hyperactivity have been linked to several types of human cancers [Senderowica, A. M., and Sausville, E. A., J. Nat. Acad. Sci., U.S.A. 96, 376-387 (2000); Garrett, M. D., Current Opin. Genetics Devel., 9, 104 (1999); Webster, K. R., Exp. Opin. Invest. Drugs, 7, 865-887 (1998); Jiang, Proc. Natl. Acad. Sci. USA 90:9026-9030 (1993); Wang, Nature 343:555-557 (1990)]. More recently, endogenous, highly specific protein inhibitors of CDKs were found to have a major affect on cellular proliferation [Sherr, C. J., Roberts, J. M. Genes Dev. 13, 1501-1512 (1999); Kamb et al, Science 264:436-440; Beach Nature 336:701-704 (1993)]. These inhibitors include p16 (an inhibitor of CDK4/cyclin D1), p21 (a general CDK inhibitor) and p27 (an inhibitor of CDK2/cyclin E). A recent crystal structure of p27 bound to CDK2/cyclin A showed how these proteins effectively inhibit the kinase activity through multiple interactions with the CDK complex [Pavletich, Nature 382:325-331 (1996)]. These proteins help to regulate the cell cycle through specific interactions with their corresponding CDK complexes. Cells deficient in these inhibitors are prone to unregulated growth and tumor formation.
In addition to treating human cancers, CDK inhibitors could be useful in the treatment of other cell proliferative disorders such as familial adenomatosis polyposis, psoriasis, neuro-fibromatosis, fungal infections, endotoxic shock, vescular smooth cell proliferation associated with atherosclerosis, pulmonary fibrosis, glomerulonephritis, and post-surgical stenosis and testenosis [U.S. Pat. No. 6,114,365].
CDKs are also know to play an important role in apoptosis. Therefore, CDK inhibitors, could be useful in the prevention of AIDS development in HIV-infected patients; inflammatory bowel disease, and diabetes mellitus, dementia, Parkinson's disease, amyotrophic lateral sclerosis, retinitis pigmentosa, spinal muscular atrophy and cerebellar degeneration; myelodysplastic syndromes, aplastic anemia, ischemic injury associated with myocardial infarctions, stroke and reperfusion injury, arrhythmia, atherosclerosis, toxin-induced or alcohol related liver diseases, hematological diseases, for example, chronice anemia; degenerative diseases of the musculoskeletal system, for example, osteoporosis, aspirin-sensentive rhinosinusitis, cystic fibrosis, multiple sclerosis, kidney diseases and pain [U.S. Pat. No. 6,107,305].
Also, it has been discovered that some CDK inhibitors can be used in combination therapy with some other anticancer agents. For example, the cytotoxic activity of the CDK inhibitor, Flavopiridol, has been used with other anticancer agents in cancer combination therapy [Cancer Research 57:3375 (1997)].
In addition, a recent report showed that CDK5 is involved in the phosphorylation of tau protein, and therefore, CDK inhibitors may be useful in the treatment of Alzheimer's disease [J. Biochem., 117: 741-749 (1995)].
This increasing body of evidence has led to intense discovery efforts to search for small molecule inhibitors of the CDK family and their associated regulatory molecules (cyclins) as an approach to cancer chemotherapy [Sausville, E. A., Trends in Molecular Medicine 8(4), S32-S37 (2002); Malumbres, M, and Barbacid, M. Nat. Rev. Cancer 1, 222-231 (2001)].
More than 50 small molecule inhibitors of cyclin-dependent kinases have been identified. These CDK inhibitors all target the ATP-binding pocket of the catalytic site of the kinases. The effects of CDK inhibitors on the cell cycle and their potential value for the treatment of cancer, alopecia, neurodegenerative disorders (e.g. Alzheimer's disease, amyotrophic lateral sclerosis and stroke), cardiovascular disorders (e.g. atherosclerosis and restenosis), glomerulonephritis, viral infections (e.g. HCMV, HIV and HSV) and parasitic protozoa (Plasmodium sp. and Leishmania sp.) has been extensively studied [Knockaert, M. Greengard, P. Meijer L., Trends in Pharmacological Sciences 23 (9), 417-425 (2002); Malumbres, M, and Barbacid, M. Nat. Rev. Cancer 1, 222-231 (2001), Sielecki, T. M. J. Med. Chem. 43, 1-18 (2000)]. Three properties make CDK inhibitors attractive as potential anti-tumor agents. First, they are potent anti-proliferative agents, arresting cells in G1 [Soni, R. J. Natl. Cancer Inst. 21, 436-446 (2001)] or G2/M [Damiens, E. et al. Oncogene 20, 3786-3797 (2001)]. Second, they trigger apoptosis, alone or in combination with other treatments [Edamatsu, H. et al. Oncogene 19, 3059-3068 (2000)]. Third, in some instances, inhibition of CDKs contributes to cell differentiation [Matushansky, I. Et al. Proc. Natl. Acad. Sci. U.S.A. 97, 14317-14322 (2000)].
Despite the significant research efforts and resources which have been directed towards the development of novel anticancer agents and improved methods for treating cancer there remains a need in the art for novel compounds, compositions, and methods that are useful for treating cancer with improved therapeutic indices.

Abstract: where G1, G2, G3, G4, A1, A2, Y1, Y2, L1, Z, e and f are defined herein, or a pharmaceutically acceptable salt thereof, which are useful for treating or preventing cancer.