Apparatus and methods for enabling a user to interact with movable control elements of a graphical user interface (GUI) by moving a hand 12, or one or more fingers 16, 18 of said hand 12 relative to a thumb 14 on said hand, in which the apparatus comprises: a camera adapted to be releasably mounted to said hand, and a processor adapted to receive signals representative of successive images captured by the camera during movement of said hand, or movement of the finger(s) of said hand relative to said thumb, to analyze said successive images to determine said movement, and to communicate with the GUI to control elements thereof in accordance with the determined movement.

Since the first commercially successful use of a GUI in the 1980s, the GUI has been adopted by most developers of computer operating systems and application software, and it is now used on or available for virtually all types of computers; GUIs are also employed, in differing forms, in a wide range of other products, such as industrial controls, televisions and other monitors, electronic games and toys. Typically, a user interacts with a GUI using a mouse, track ball, touch pad or the like, often in concert with a keyboard; interaction is also possible with a keyboard alone, however users find it faster, easier and often more intuitive to use a mouse, etc. rather than just a keyboard. The use of such device combines with the advantageously intuitive characteristics of the GUI to provide users with a simple way to interact with computers and the like, such that the use of such devices as mice, touchpads, trackballs, joysticks, etc. to provide a primary means of interaction with a GUI is almost ubiquitous today (with a keyboard, console or other control providing an important, secondary means of interaction).
The rapid and continuing development of electronic devices and software, and the massive increases in processing power and memory now available to the general population have begun to highlight limitations and disadvantages of typical GUI interface or interaction devices such as mice, touchpads, track balls and the like. Mice need to be positively held by a user, and need a preferably flat surface upon which they can be moved in order to effect many of their functions. Because they are designed to be moved over a flat surface, mice can only effectively provide control over a two-dimensional GUI (which limits their use in applications such as computer aided design or games, which are very often three-dimensional, or with the emerging three-dimensional displays and TVs); moreover, the range of control movement offered by a mouse is restricted by the size of the surface upon which the mouse sits (which is one reason why mice often incorporate a "scroll wheel"). Such limitations and disadvantages apply, to a greater or lesser extent, to similar GUI interface devices such as touchpads, trackballs and the like.
In computer gaming, manufacturers have attempted to address interaction with the GUI (or game), in different ways: initially, control/interaction was carried out through specialised consoles, which incorporated many of the functions of mice, etc., but in a different format. These consoles began to be equipped with accelerometers, or other means which were able to recognise motion of the console (such as by the console having an infra-red transmitter, with movement of the console being monitored by remote infra-red sensors) and to translate this recognised movement into instructions for player-game interaction. Consoles designed to be held in two hands were replaced by smaller consoles which a user could hold in a single hand, allowing a user to play electronic games such as tennis or golf. A later and ongoing advance in the gaming field is to provide a way of registering movement or gestures made by the user (as opposed to something held or carried by the user) and translating these into instructions for the game; such an approach is now also being adopted by the manufacturers of other equipment, such as the so-called "smart" TVs, where it is becoming more common to enable a user to interact with the GUI by way of gestures alone. Such an approach is attractive, because it facilitates a very intuitive interaction, however it has a number of disadvantages. It requires sophisticated and sensitive sensors in order first to capture the movement of a user, then it requires significant amounts of processing capability to analyse the movement and to translate it into the instructions or controls which are intended by the user and also causing significant latency; as a result, such systems today are reasonably reliable in interpreting gross movements and converting these into controls to the GUI, but less so where smaller or more precise and/or rapid movements or gestures are concerned. They often have a limited field of view, limiting the range of movement of the user, or require some way of tracking the user, they risk being confused by movements other than by those of the user, and because they are most reliable with gross movements they are inherently slow to respond and limited in the range of different functionalities they can currently provide, because the range of gross movements available to a user is limited.


A binder for preparing a positive electrode of a lithium secondary battery, and a method for preparing a positive electrode using the same. The binder includes two or more different lithium-substituted polyacrylic acids with different molecular weights. The lithium-substituted polyacrylic acids include two different lithium-substituted polyacrylic acids differing in weight average molecular weight by 500,000 or more from each other.

As application areas of secondary batteries expand to electric vehicles (EV), energy storage systems (ESS) or the like, lithium-ion secondary batteries having relative low weight to energy storage density (˜250 Wh/kg) have limits in applications for such products. On the other hand, a lithium-sulfur battery among lithium secondary batteries is capable of accomplishing high weight to energy storage density (˜2,600 Wh/kg) theoretically, and thereby has received attention as a next generation secondary battery technology.
A lithium-sulfur battery means a battery system using a sulfur series material having a sulfur-sulfur (S—S) bond as a positive electrode active material, and lithium metal as a negative electrode active material. Sulfur, a main material of the positive electrode active material, has advantages of being very abundant in resources globally, having no toxicity and having a low atomic weight.
During discharge of a lithium-sulfur battery, lithium, a negative electrode active material, is oxidized while releasing electrons and being ionized, and a sulfur series material, a positive electrode active material, is reduced by receiving the electrons. Herein, the oxidation reaction of lithium is a process in which lithium metal releases electrons and changes into a lithium cation form. In addition, the reduction reaction of sulfur is a process in which an S—S bond receives two electrons and changes into a sulfur anion form. The lithium cation produced through the oxidation reaction of lithium is transferred to a positive electrode through an electrolyte, and forms a salt by bonding with the sulfur anion produced through the reduction reaction of sulfur. Specifically, sulfur before discharge has a cyclic S8 structure, and it changes to lithium polysulfide (LiSx) through a reduction reaction. When the lithium polysulfide is fully reduced, lithium sulfide (Li2S) is produced.
Despite the lithium-sulfur battery having an advantage of high energy storage density, many problems are present in actual applications. Specifically, a problem of instability of lithium metal used as a negative electrode, a problem of low conductivity of a positive electrode, a problem of sulfur series material sublimation when preparing an electrode, a problem of sulfur series material loss during repeated charge and discharge, and the like, may be present.
Existing positive electrodes for a lithium-sulfur battery have been prepared using styrene butadiene rubber/carboxymethyl cellulose (SBR/CMC) as a binder. However, dispersibility was reduced when a thickener such as CMC was not present in the binder, and a problem of a long-term stability decrease caused by positive electrode deterioration as charge and discharge proceed has been identified.
Accordingly, binders capable of improving lithium-sulfur battery performance still have been required in the art.

Binder for manufacturing positive electrode of lithium secondary battery and method for manufacturing positive electrode by using same


This specification describes a completely new and different concept, method, mechanization, apparatus configuration and sequencing procedure for obtaining commercially useful energy and power, namely: the use of a gyroscope to generate output power from the earth's stored inertial rotational (flywheel) energy by fixing the housing of the gyroscope to the earth and using the rotation of the earth relative to the gyroscope's spatially stable rotor/gimbal assembly to rotate the input shaft(s) of a power transducer(s). This provides a power-output stroke over a near 180.degree. precession excursion (near pole-to-pole alignment) until the rotor spin axis of the gyroscope is nearly aligned with the polar axis of the earth. Re-precession torques are then created on the inner gimbal within the gyroscope to re-orient the rotor spin axis vector direction by re-precessing the rotor spin by somewhat more than 180.degree. to the near polar alignment direction required for the resumption of power output in a manner requiring essentially zero energy and power to provide relatively continuous, but intermittent, power output generation along with a multi-unit assembly for producing continuous uninterrupted power output.. An "Advanced Concept" based on rate gyro characteristics showing promise of enhanced performance is included.

The angular space-stability characteristic of gyroscopes relative to the earth's rotational movement is known. The use of precession of the spin axis of a gyroscope rotor, where the gyroscope housing is fixed relative to a host vehicle and the rotor/gimbal assembly is fixed to motionless angular space in order to generate navigational signals, is also well known.
However, it is not known how, without exerting a great amount of excitation (input) power, to effectively use a gyroscope in order to obtain significantly useful net output power via gyroscopic-rotor-precession-generated reaction torque from the earth's vast storehouse of rotational inertial energy on a continuous basis.
One reason for this is that the gyroscope rotor/gimbal assembly is stable (unmovable and unperturbable) when the rotor spin axis is in alignment with (parallel to) the earth's polar axis--a conceptual "catastrophe" when first encountered during preliminary concept analysis.
At this orientation of the rotor spin axis, the precession resulting from the application of precession torque of the rotor spin axis stops, and no further gyro reaction torque can be produced by the gyro rotor/gimbal assembly to oppose the earth's angular rotational velocity imposed torque as expressed by the gyro performance equation T.sub.p =K.sub.p.OMEGA..sub.p. When the rotor precession rate .OMEGA..sub.p =0, T.sub.p =0.
Thus, there is no further access to the earth's stored rotational energy to generate useful output power--unless corrective measures are applied.
These corrective measures are the primary subject of this application for letters patent.
The conceptual theoretical solution to the "polar axis alignment catastrophe" would be to re-orient the rotor spin vector, somehow, back to the starting orientation of the power output precession "stroke",--hopefully, in a manner requiring a negligible amount of input or excitation energy compared to the energy extracted during the energy/power-output phase.
Some fairly obvious means exist for physically re-orientating the rotor spin vector direction of the gyroscope from the undesired stable polar axis alignment orientation to one wherein the spin vector direction of the rotor is such that it will again traverse a slightly less than 180.degree. precession angular displacement toward the stable polar axis alignment orientation, during which angular excursion power may be extracted from the earth's rotational energy. Such obvious and simple re-orienting means generally require substantial energy expenditure. Thus, they are not practical from a power generation efficiency standpoint.
For example, this can be accomplished by simply reversing the direction of rotor spin. But the energy required to "brute force" reverse the rotor spin direction far exceeds the energy obtained during the intervening power output phase of the operational cycles. Even if the net energy output were positive, the time required for rotor de-re-spin is excessive, making the overall time-averaged power efficiency unattractive.
Re-precession torque applied around the torque power output (polar) axis by torque reversal, while properly re-orienting the rotor spin axis, is counterproductive in that, since the direction of the earth's rotation remains the same and the direction of the torque is reversed, the direction of energy flow is reversed. Thus, this procedure merely results in returning the energy to the earth's rotation that was extracted during the preceding power-output phase--plus a "tax" for the always attendant parasitic inefficiency factors.
So, to tap the physical rotational energy of the earth, a new, effective, and efficient method, apparatus, and mechanization is needed to produce useful power-output from the earth's vast storehouse of inertial rotational energy.
This need can be satisfied by employing the angular spatial stability and torque production characteristics of a precessing gyroscope rotor. By applying the torque production capability of a precessing gyro rotor as a deceleration torque on the earth's angular momentum, the earth's storehouse of rotational inertial energy may be accessed by interposing a power transducer between the motionless (in spatial coordinates) gyro outer-gimbal and the earthly inertial rotation.
To alter the earth's spatial parameters has been believed unattainable ever since Archimedes, 2500 years ago, requested ". . . a pole (lever) long enough and a platform . . . [in order to] move the earth." That platform and lever (torque) are now readily available, and have been, but unrecognized as such, for over a century. It's a gyroscope rotor/gimbal assembly.
The practical extraction of this readily available supply of enormous energy needs to be accomplished in a manner that does not require the expenditure of a large amount of excitation power, or include complex electrical or mechanical operations on the gyroscope and rotor, or that is otherwise impractical.
The basic elements of this invention were disclosed in Disclosure Document No. 220,425, filed Feb. 21, 1989, resulting from the first year of research on the concept.
This application for letters patent is based upon continuing R&D conducted over the ensuing elapsed time to "prove out" this new and different conceptual configuration of a gyroscope designed to extract power and energy from the inertial energy of a rotating planetrary body, such as the earth, consisting of the following events and milestones:
(1) The determination that theoretically there is an enormous quantity of pure unpolluting energy available from the energy contained in the rotational inertia of the earth--enough to supply all the energy and power requirements of the world's civilization, currently being met by oil production, for the next 100,000,000 years with only minimal and readily acceptable reduction in the daily rotation rate of the earth at the end of that period: approximately 30%, a 34 hour day--at the present rate of oil energy consumption.
(2) The determination that this supply of pure physical mechanical energy may theoretically be tapped by interposing the spatial angular stability of a gyro rotor/gimbal assembly, plus a transducer to convert the gyro/earth mechanical energy and power into the desired type for ready utilization, e.g., electrical, hydraulic, or other, between the onrushing inertial rotational energy of the earth and the referenced spatial stability--more accurately, rigidity (as in "brick wall" or Hoover Dam)--of a gyro rotor and gimbal assembly.
(3) Theoretical analysis showing that by applying gyro precession reaction torque induced by the spatially rotating earth to a power transducer, or generator, shaft, the housing or framework of both gyro and generator being anchored to the earth, there is a direct correlation between: (a) the theoretically/mathematically determined reduction in rotational energy of the earth and (b) the amount of theoretically/mathematically determined earth-derived mechanical energy delivered to the input shaft of the power transducer (electrical power generator).
(4) Realization that re-precessional re-orientation performed by applying re-precession torque around the polar axis, while accomplishing the desired re-orientation of the spin vector, would be counterproductive. Since the applied torque would be opposite to that employed during the power output phase and the earth rotation direction remains the same, the direction of energy flow would be reversed. Thus the rotational energy of the earth would be increased at the expense of some other source of energy.
(5) Realization that by definition there is no rotation of the earth around an axis in, or parallel to, the earth's equatorial plane. Therefore, applying re-precession torque around an axis in, or parallel to, the earth's equatorial plane, theoretically would require only infinitessimal energy or power in comparison to that obtainable by means of polar axis precessional torque.
(6) Since there are three orthogonal axes in the gyro/earth coordinate system and torque around one axis creates precession around an axis normal to that axis, for any precession (motion) axis there are two precession torque axes that can cause precession around that same (motion) axis. Since the desired re-precession motion axis is in the equatorial plane, in order to re-precess or re-orient around the same equatorial axis as the power-output/polar-axis torque created precession motion, and the polar axis cannot be used for re-precession, the remaining, or third, axis lies in the equatorial plane normal to the desired re-precession action axis. Re-precession torque around this axis should theoretically require essentially no energy or power from the re-precession torque excitation source.
(7) Mathematical analysis of re-precession torque applied around an equatorial axis confirmed that the energy and power required to re-precess with torque around, or parallel to, an equatorial plane axis was indeed a theoretical, mathematical, and practical infinitesimal, many orders of magnitude below that obtained in the power output phase.
(8) Laboratory experimentation confirmed the validity of the theoretical/mathematical gyro/earth-rotational power output performance equations and that a high degree of correlation could be demonstrated between theory and practice.


Disclosed is a technique for reducing noise superimposed on an output signal while keeping loop gain constant without increasing the circuit scale and without changing the transfer function of the amplifier apparatus (frequency characteristics of gain and phase). According to the technique, there are included a power-supply voltage control unit 7 for detecting the amplitude level S9 of an input audio signal S1 and outputting power with a voltage value indicated by target set voltage value information Vs corresponding to this amplitude level S9, and a PWM modulation unit 2 including a PWM converter 23 for converting the pulse width of the input audio signal S1 and a correction unit for correcting the signal modulated by the PWM converter 23. The PWM modulation unit 2 corrects the pulse width of a PWM signal S5 modulated by the PWM converter 23 so that the correction unit will cancel out a change in amplification gain of a power amplification unit 4 according to the target set voltage value information Vs.

Conventionally, there has been a technique for using a voltage-variable power supply as the power supply of an amplifier apparatus to make it follow the level of an input audio signal to the amplifier apparatus in order to increase or decrease the value of power supply voltage to be supplied to a power amplifier stage, thereby reducing noise superimposed on an output signal and improving the power efficiency at the power amplifier stage. In this case, the power supply voltage to be supplied to the power amplifier stage is made to follow the level of the input audio signal. Therefore, when the level of the input audio signal is low, since the power supply voltage at the power amplifier stage can be reduced to a voltage value with an amplitude that does not distort the amplified signal, noise superimposed on the output signal of the amplifier apparatus can be reduced and power efficiency at the power amplifier stage can be improved.
However, when the power supply voltage to be supplied to the power amplifier stage is lowered, the amplification gain of the amplifier apparatus due to voltage amplification at the power amplifier stage also drops. To prevent this, there is known a technique in which a constant voltage control circuit controls the voltage of a power amplification unit in proportion to the amplitude level of an input analog signal, and an attenuator provided in a feedback circuit cooperates with the constant voltage control circuit to adjust a feedback amount to be kept constant according to the voltage level of a constant voltage (for example, see Patent document 1 cited below). There is also known another conventional technique in which the gain of gain adjustment means provided in a feedback loop is adjusted according to an increase or decrease in voltage value to be supplied to a switching amplifier in a signal reproducing apparatus (for example, see Patent document 2 cited below).

PANASONIC INTELLECTUAL PROPERTY MANAGEMENT CO., LTD.


A thermal monitor diode is provided that comprises a silicon thin film on an insulator mounted on a silicon substrate. An opening extends through the silicon thin film and through the insulator and partially into the silicon substrate and terminates at an end wall. A conductive material is disposed in the opening and extends to the end wall. The substrate has a P/N junction formed therein adjacent the end wall, and an insulating spacer material surrounds the conductive material and is sufficiently thin to allow temperature excursions in the silicon thin film to pass therethrough. The invention also contemplates a method of forming the diode.

In thin film technology for the formation of integrated circuit (I/C) chips, it is often necessary to monitor and/or control the internal chip, sub circuit, or voltage island temperature. According to the prior art, a typical technique for sensing such temperatures is the use of a thermocouple. However, in many instances, the accuracy of the thermocouple is inadequate to perform the real time automatic control of the I/C chip or chips necessary or desirable. It is, therefore, desirable to provide a rapid response time temperature sensor and/or controller in thin film I/C chip circuit technology.

Autonomic thermal monitor and controller for thin film devices

International Business Machines Corporation


A rolling walker comprises a frame, a seat supported by that frame, and a backrest supported by the frame. By one approach the backrest is configured to selectively move between a first position that provides back support for a person sitting in the seat facing in a forward direction and a second position that provides back support for a person sitting in the seat facing in a rearward direction. If desired, this backrest can be comprised of a material (such as a memory foam material) that biases the backrest towards that first position when the backrest is in the first position and that biases the backrest towards the second position when the backrest is in the second position.

Wheelchairs are typically designed to transport a sitting person and so-called companion chairs are a lighter-duty mechanism having a similar operating purpose. Accordingly, both wheelchairs and companion chairs typically have leg riggings to support the transportee's lower appendages above the ground. Rollators are a walking aid and hence lack such leg riggings. That said, some rollators include a seat. This seat provides the user with a place to sit when that need arises (for example, when the user needs a break from standing or walking).
The basic design for a rollator is well established; a frame having four ground-contacting wheels and a pair of handles that the user can grip when walking with the aid of the rollator. Unfortunately, these deceptively simple design concepts are not always implemented in a fashion that well suits the needs of the expected user population. Persons who seek walking assistance can also present a variety of other maladies, infirmities, and conditions that can, in practice, interfere with their successful use of the rollator. Examples include, but are not limited to, reduced dexterity or upper-body strength, limited visual acuity, and reduced cognitive capabilities.
The rollator user population also represents a wide variety of usage patterns, lifestyles, and operating environments. Some users, for example, may only utilize their rollator within a fairly limited and constrained application setting while other users may need to frequently transport their rollators in a vehicle and more aggressively use their rollators in a variety of application settings.
Elements in the figures are illustrated for simplicity and clarity though are drawn to scale. Also, common but well-understood elements that are useful or necessary in a commercially feasible embodiment are often not depicted in order to facilitate a less obstructed view of these various embodiments. The terms and expressions used herein have the ordinary technical meaning as is accorded to such terms and expressions by persons skilled in the technical field as set forth above except where different specific meanings have otherwise been set forth herein.


A position detecting system includes an RFID terminal configured to transmit an identification information signal by radio; a plurality of receiving devices configured to receive the identification information signal; and a position detecting server configured to receive a reception notification of the identification information signal from each of the receiving device and configured to estimate a position of the RFID terminal. The RFID terminal includes a part configured to perform radio transmission of a plurality of identification information signals having different radio strengths. The position detecting server includes a part configured to integrate the reception notifications of the identification information signals having different radio strengths and configured to estimate the position of the RFID terminal.

1. Field of the Invention
The present invention generally relates to position detecting systems and active RFID (Radio Frequency Identification) terminals. More particularly, the present invention relates to a position detecting system for position detecting by using an active RFID terminal and the RFID terminal.
2. Description of the Related Art
FIG. 1 is a structural view of an example of a related art position detecting system. As shown in FIG. 1, plural receivers 12a through 12j are fixed within a position detecting area 10. In addition, an active RFID terminal 14 moves in the position detecting area 10.
The active RFID terminal 14 sends a identification information signal by itself within a designate time interval. A reaching area (receivable area) of the identification information signal is shown by a circle 15. Because of this, the receivers 12a, 12b, 12c, and 12d situated in the circle 15 receive the identification information signal from the active RFID terminal 14.
When each of the receivers 12a through 12j receives the identification information signal, ID information demodulated from the identification information signal is sent to a server 16 via, for example, a LAN (Local Area Network).
When the server 16 receives the notification of the ID information from the receivers 12a, 12b, 12c, and 12d, which have received the identification information signal from the active RFID terminal 14, the server 16 calculates center of gravity 17 of positions of the receivers 12a, 12b, 12c, and 12d so as to estimate the center of gravity 17 as a present position of the active RFID terminal 14.
FIG. 2 is a block diagram of the RFID terminal 14 shown in FIG. 1. As show in FIG. 2, the ID information for identifying the RFID terminal is stored in an ID-ROM part 20. A transmission control part 21 codes the ID information read from the ID-ROM part 20 and supplies coded ID information to the RF part 22 with a designated time interval. The RF part 22 modulates the coded ID information and transmits it as an identification information signal of a specific frequency from an antenna 23.
FIG. 3 is a structural view of another example of the related art position detecting system. In FIG. 3, parts that are the same as the parts shown in FIG. 1 are given the same reference numerals, and explanation thereof is omitted.
As shown in FIG. 3, when a server 18 receives the ID information from the receivers 12a, 12b, 12c, and 12d receiving the identification information signal from the active RFID terminal 14, the server 18 weights each of positions of the receivers 12a, 12b, 12c, and 12d by using larger weights as the number of identification information signals received by the receivers 12a, 12b, 12c, and 12d is larger. Then the server 18 calculates the center of gravity 19 so as to estimate the center of gravity 19 as a present position of the active RFID terminal 14.
Here, the number of receptions of the identification information signals by the receiver 12c near the active RFID terminal 14 is large. The number of receiving the identification information signal by the receiver 12a far from the active RFID terminal 14 is small. The present position of the active RFID terminal 14 is estimated by being biased to a side of the receiver 12c where the number of received signal is large by weighting the number of received signals so that the position estimation precision of the active RFID terminal 14 is improved.
Patent Publication 1 mentioned below describes that an electric wave transmitted from the RFID is received by plural receivers so that the position of the RFID is specified by positions of the plural receivers.
Patent Publication 2 mentioned below describes that an electric wave transmitted from the portable terminal is received by plural receivers so that the position of the portable terminal is specified based on each of received levels at which plural receivers receive.
Patent Publication 3 mentioned below describes that an electric wave transmitted from the RFID is received by plural receivers and a position corresponding to the receiver closest to the RFID is specified as the position of the RFID based on the received radio strength of plural receivers.
Patent Publication 4 mentioned below describes that an electric wave transmitted from the portable terminal is received by plural local control stations and the position of the portable terminal is specified by position of the local control station or received radio strength.    [Patent Document 1] Japanese Laid-Open Patent Application Publication No. 10-124781    [Patent Document 2] Japanese Laid-Open Patent Application Publication No. 8-307917    [Patent Document 3] Japanese Laid-Open Patent Application Publication No. 11-224400    [Patent Document 4] Japanese Laid-Open Patent Application Publication No. 11-169353
However, in the above-discussed related art cases, the position of the RFID is specified based on the position of the receiver which receives the identification information signal or by weighting the number of receptions of the identification information signal to the position of the receiver. Therefore, there is limitation of precision for specifying the position of the RFID. In addition, in an area where the electric wave reaches, depending on the difference of frequency of the electric wave or reflection, there may be a position where the electric wave is not easily received by the receiver. This may cause degradation of precision for specifying the position of the RFID.

Position detecting system and position detecting server


A sensor unit for use in a surveillance system is provided. The sensor unit interfaces with a sensor device to receive sensor data from the sensor device. The sensor unit is configured to generate an event record based upon received sensor data. The event record is of a predetermined format and is output to a data unit.

In typical surveillance systems one or more sensor devices are used to capture/respond to particular conditions, changes or occurrences. Signals from the sensors devices are provided to a monitoring unit to indicate the particular condition, change or occurrence. When the monitoring unit receives these signals, an alert may be generated to advise of the detected condition/change. In the case where the sensor is, for example, an imager, such as a video camera, the signal from the sensor may be presented for display in real time on a display device and/or recorded to a recycling recording device, such as a linear video tape recorder.
Other than recording information to a recording device, no data concerning the change, condition or occurrence is collected or otherwise generated. Thus, the ability to analyze/evaluate the change, condition or occurrence is very limited. Further, the ability to take appropriate or necessary action for a given situation is also limited. In general, unless a party is available at the monitoring station to receive the alert or view information on a display; very little analysis of the condition, change or occurrence is possible. Further, the meaning or relation of the condition, change or occurrence with respect to past or future conditions, changes or occurrences is not considered or otherwise taken into account by the typical surveillance system.
Thus, a heretofore unaddressed need exists in the industry to address the aforementioned deficiencies and inadequacies.


An image forming apparatus includes: an image carrier on which an electrostatic latent image is to be formed; a rotary development unit having a plurality of developer carriers that carry a developer used for developing the electrostatic latent image, the developer including toner, the plurality of the developer carriers being sequentially moved to a development position opposing to the image carrier, to switch a development color; and a toner concentration detection unit that detects the concentration of toner of the developer carried by one of the developer carriers located at the development position.

1. Field of the Invention
The present invention relates to an image forming apparatus which forms an image by means of electrophotography.
2. Description of the Related Art
Among pieces of image forming apparatus which form an image by means of electrophotography, such as a full-color copier, a printer, and a multifunction machine, one type is equipped with a rotary development unit having four development devices assigned to respective Y (yellow), H (magenta), C (cyan), and K (black) colors. In the image forming apparatus of this type, each of the development devices of the rotary development unit is provided with one development roller.
Rendering the density of an image stable has recently become important for enhancing image quality. Therefore, when a two-component developer formed from toner and a carrier is used, the concentration of toner in the developer (hereinafter simply called "toner concentration" or a "toner mixing ratio") must be appropriately controlled. Methods for detecting a toner concentration include, in addition to a method for detecting the concentration of toner in a developer agitated in a development device, a method for indirectly ascertaining the concentration of toner in the development device by means of detecting the concentration of toner in the developer carried by a development roller.
An image forming apparatus equipped with a toner concentration sensor (TC sensor) for detecting the concentration of toner on a development roller has become known as a related-art image forming apparatus having a rotary development unit. In such a conventional image forming apparatus, the concentration of toner on the development roller is detected by a toner concentration sensor. On the basis of a result of detection, supply of toner to the development apparatus from the toner cartridge is controlled. Such a system for controlling toner supply using a toner concentration sensor is also called an Auto Toner Concentration (ATC) system.
However, the related-art image forming apparatus adopts a construction for detecting the concentration of toner on a development roller through use of a toner concentration sensor and at a position differing from a position where an electrostatic latent image is actually developed (hereinafter described as a "development position"); e.g., a position displaced from the development position by 90° or 180° in the rotational direction of the rotary development unit. The apparatus suffers the following drawbacks.
Specifically, when toner supply is controlled by means of the ATC system, the concentration of toner on the development roller must first be detected at a position differing from the development position. Subsequently, the development roller must be replenished with toner after having been moved to the development position. Therefore, the image forming apparatus equipped with a rotary development unit cannot control the toner supply in real time.

Image forming apparatus with toner concentration detection unit


A plurality of piezoelectric elements are arrayed two-dimensionally. A plurality of electrodes are respectively formed on the plurality of piezoelectric elements. A plurality of non-conductive members have columnar shape and are arranged on the plurality of electrodes. A plurality of internal metal layers are respectively provided for the plurality of non-conductive members. The internal metal layers reach from arrangement surfaces of the non-conductive members to other surfaces of the non-conductive members. The arrangement surfaces are opposite to the other surfaces.

1. Field of the Invention
The present invention relates to an ultrasonic probe and an ultrasonic diagnosis apparatus having a two-dimensional array structure, and an ultrasonic probe manufacturing method.
2. Description of the Related Art
An ultrasonic probe having a one-dimensional array structure is available. A transducer unit included in this one-dimensional array ultrasonic probe has a plurality of transducers arrayed in a line. In general, electrodes on the upper and lower surfaces of the transducer unit are extracted from an end of the transducer unit. Various contrivances have been made to extract upper surface electrodes. For example, there is available a technique of electrically extracting upper surface electrodes from the lower surface of a transducer unit via an FPC (Flexible Printed Circuit board) by plating a side surface of the transducer unit to render the upper and lower surfaces conductive. The signals extracted by the FPC are transmitted to a transmission/reception circuit via a probe cable.
In general, the acoustic impedance of polyimide used as a base material for an FPC is about 3 MRayl. The acoustic impedance of a transducer unit is equal to or more than 30 MRayl. For this reason, when the FPC is directly joined to the transducer unit, an acoustic mismatch occurs. In order to reduce this acoustic mismatch, an acoustic matching layer having an acoustic impedance between 3 MRayl and 30 MRayl is used. This acoustic matching layer is placed on the upper surface of the transducer unit, and the FPC is placed on the upper surface of the placed acoustic matching layer. Upper surface electrodes are electrically extracted via this FPC.
In the case of specifications with three acoustic matching layers added to a transducer unit, the first acoustic matching layer has the best acoustic impedance of about 9 to 15 MRayl. A material having such an acoustic impedance is a ceramic material containing mica as a main component. This ceramic material is known as a machinable ceramic material. This material has non-conductivity. There is a technique uses a method of plating all the surfaces of the first acoustic matching layer using this non-conductive material and electrically extracting upper surface electrodes formed out of a piezoelectric element to the upper surface of the acoustic matching layer.
In a three-layer specification two-dimensional array ultrasonic probe, a multilayer structure comprising a plate-like piezoelectric member, a first acoustic matching layer member, and a second acoustic matching layer member is cut in a lattice form. With this cutting, each acoustic matching layer is divided into a plurality of acoustic matching elements arrayed two-dimensionally. In the above method of extracting upper surface electrodes by plating the surrounding portion, the upper and lower surfaces of acoustic matching elements other than those located outside the first acoustic matching layer are not rendered conductive.
As another method of electrically extracting upper surface electrodes to the upper surface of an acoustic matching layer, a method of attaching a conductive pattern to a side surface of an acoustic matching layer has been proposed. In this method, however, pattern attachment processing needs to be performed for each column. This increases the number of steps, resulting in an increase in cost.

Ultrasonic probe, ultrasonic diagnosis apparatus, and ultrasonic probe manufacturing method


A display panel including a display part including a plurality of sub-pixels configured to display a plurality of colors, and a plurality of data lines connected with the sub-pixels; a first test part configured to supply a test signal to (2K−1)th data lines (‘K’ is an integer above 0) by each color for the sub-pixels among the plurality of data lines; and a second test part configured to supply a test signal to 2Kth data lines by each color for the sub-pixels among the plurality of data lines when the first test part supplies the test signal. Further, a polarity of the test signal supplied by the second test part is opposite to a polarity of the test signal supplied by the first test part.

Field of the Invention
Embodiments of the present invention relate to a method for testing whether or not a display panel is normally operated.
Discussion of the Related Art
Display devices such as a liquid crystal display (LCD), organic light-emitting diodes (OLED), plasma display panel (PDP) and electrophoretic display (EPD) may be manufactured by various processes including a process of testing lighting for a display panel. In more detail, an image is displayed on the display panel included in the display device. The process of testing lighting can be performed by supplying a predetermined test signal to the display panel, and testing whether or not the display panel is normally operated according to the predetermined test signal.
For example, FIG. 1 is a schematic view illustrating a related art display panel, and FIG. 2 is a schematic view illustrating the display panel when a process of testing lighting is performed.
Referring to FIG. 1, the related art display panel 10 includes a display part 11 and a testing part 12. The display part 11 includes a plurality of sub-pixels SP defined by a plurality of gate lines GL1 to GLn and a plurality of data lines DL1 to DLm, in which the gate line and the data line cross each other. The plurality of data lines DL1 to DLm are also connected to the testing part 12.
Further, the testing part 12 alternately supplies a positive polarity (+) test signal and a negative polarity (−) test signal with respect to a common voltage Vcom to the data lines DL1 to DLm. That is, the testing part 12 transmit the positive polarity (+) test signal and the negative polarity (−) test signal, which are supplied from a testing apparatus, to the data lines DL1 to DLm, whereby the positive polarity (+) test signal and the negative polarity (−) test signal are alternately supplied to the data lines DL1 to DLm. The data lines DL1 to DLm are also connected to the testing part 12 through one connection line.
As the testing part 12 alternately supplies the positive polarity (+) test signal and the negative polarity (−) test signal to the data lines DL1 to DLm, the sub-pixels SP included in the display panel 10 according to the related art are driven and emit light in units of a frame, to thereby perform a lighting test. For example, all the sub-pixels SP of the display panel 10 according to the related art are driven according to the positive polarity (+) test signal during a first frame, and then all the sub-pixels SP are driven according to the negative polarity (−) test signal during a second frame.
According to these repetitive driving operations, the lighting test is performed. That is, a frame inversion method is performed by repetitively inverting a voltage polarity in unit of a frame in all the sub-pixels SP of the display panel 10 according to the related art.
However, because the lighting test is performed in the frame inversion method, a flicker occurs due to a low frequency for the test signal, which causes problems in the lighting test. In order to overcome this flicker problem, a method of increasing a frequency for the test signal has been proposed. However, if the frequency is increased for the test signal, the time for applying the positive polarity (+) test signal and negative polarity (−) test signal to the data lines DL1 to DLm becomes short. Thus, as shown in FIG. 2, some of the sub-pixels SP included in the display panel 10 are driven in an incomplete charging state, whereby a lighting level in some sub-pixels SP driven in the incomplete charging state is relatively lower. Therefore, based on the testing results, even though the sub-pixels SP have a good quality, a classification the lighting level in the sub-pixels SP is marked as having a poor quality.

Display panel and method for testing display panel


An input device includes a selection circuit that is connected to a plurality of electrodes provided on a capacitive touch panel and selectively connects the plurality of electrodes, a calibration resistor connected to the selection circuit, a constant current source for supplying constant current power to the plurality of electrodes and the calibration resistor, an A/D converter connected to the selection circuit to measure capacitance of the plurality of electrodes and a measured value measured when the calibration resistor is connected, and a control unit controlling the selection circuit to make the A/D converter measure the capacitance of one of the plurality of electrodes and the measured value measured when the calibration resistor is connected, and determining a contact coordinate of the touch object on a touch panel by calculating a ratio of the capacitance to the measured value.

1. Field of the Invention
The invention relates to an input device.
2. Description of the Related Art
An input device is known that a false detection is prevented by measuring time for charging electrodes on a capacitive touch panel and time for discharging them (see, e.g., JP-A-2011-076515).
The input device has plural pairs of electrodes two-dimensionally arranged in a touch detection area, a selection circuit sequentially selecting each of the plural pairs of electrodes, a charge circuit for charging a pair of electrodes selected by the selection circuit, a charge time measuring circuit for measuring charge time from when the charge circuit starts charging the pair of electrodes until reaching a predetermined upper threshold, a discharge circuit for discharging the pair of charged electrodes and a discharge time measuring circuit for measuring discharge time from when the discharge circuit starts discharging the pair of electrodes charged to the predetermined upper threshold until reaching a predetermined lower threshold, and the input device determines, based on a ratio of the charge time to the discharge time, whether contact of a touch object in a touch area is detected or it is false detection.

KABUSHIKI KAISHA TOKAI RIKA DENKI SEISAKUSHO


In a drive apparatus for a toggle-type mold clamping mechanism of a vertical injection molding machine, a drive motor movably supported on a guide table via rollers moves horizontally following the vertical movement of a ball screw, when a mold is opened or closed, with keeping the distance between the ball screw and the output shaft of the drive motor constant by means of connecting plates. As the ball screw is rotated by the drive motor via a first pulley, a second pulley and a belt wound around the first and second pulleys, a toggle arm of each of a pair of symmetrical toggle clamping mechanisms is bent to close or open the mold.

1. Field of the Invention
This invention relates to a drive apparatus for a toggle-type mold clamping mechanism of a vertical injection molding machine.
2. Description of the Related Art:
As shown in FIG. 10 of the accompanying drawings, a typical conventional drive apparatus for a toggle-type mold clamping mechanism of a vertical injection molding machine comprises a pair of toggle mechanisms, each including a pair of links 101, 101 pivotally joined at their inner ends and pivotally connected at their outer ends to a stationary platen 102 and a rear platen 103, respectively, a nut 106 to which the inner ends of the links of one toggle mechanism are connected, a ball screw on which the nut 106 is threadedly mounted, and a motor 107 to which the inner ends of the links of the other toggle mechanism are pivotally connected and which has an output shaft secured to the ball screw. When the distance between the inner ends of the links of one toggle mechanism and those of the links of the other toggle mechanism is varied, mold halves 105 secured to the bottom surface of an upper plate 104 and the top surface of the stationary platen 102, respectively, are opened or closed.
With this conventional drive apparatus for a toggle-type mold clamping mechanism, the motor is attached in an unbalanced manner to the movable parts which move vertically simultaneously with the vertical movement of the upper and rear platens. It is impossible, therefore, to keep the mold-attaching surfaces of the stationary and upper platens parallel to each other due to the weight of the movable parts and due to a deviated load on the movable parts during closing and opening of the mold halves. Further, since the weight of the movable parts becomes increased, it is time-consuming to accelerate and decelerate the motor so that the mold halves can hardly be opened and closed at high speed.


An image sensor includes a sampling circuit, which is responsive to ramp and input signals provided by a pixel array. The sampling circuit is configured to perform a correlated double sampling (CDS) operation on a reset component of the input signal and an image component of the input signal. A reset circuit is also provided, which is configured to reset an output signal generated by the sampling circuit to a first level in response to a reset control signal. This reset control signal can be active during at least one of a first reset interval and a second reset interval. The first reset interval precedes a first comparison interval during which a first comparison operation is performed on the ramp signal and the reset component. The second reset interval precedes a second comparison interval during which a second comparison operation is performed on the ramp signal and the image component.

1. Technical Field
Example embodiments relate to correlated double sampling, and more particularly to correlated double sampling circuits, analog to digital converters and image sensors including the same.
2. Description of the Related Art
Charge coupled device (CCD) image sensors and complementary metal-oxide semiconductor (CMOS) image sensors have been used as devices for taking images. In general, analog pixel signals output from pixel arrays of the CMOS image sensors may have variations due to the difference between intrinsic characteristics of pixels, such as fixed pattern noise (FPN). To reduce the FPN, a correlated double sampling (CDS) technology may be used to operate the CMOS image sensors.

Correlated double sampling circuit, analog to digital converter and image sensor including the same


A charge pump latch circuit is provided that includes at least one first and at least one second charge pump stage interconnected by an intermediate circuit node, and a stabilization stage connected to the intermediate circuit node and to control terminals of transistors of the first and second charge pump stages. The stabilization stage includes at least one first pair and at least one second pair of first and second enable terminals receiving suitable and distinct phase signals that ensure the turn-off of the stabilization stage during overlapping periods of the phase signals. Also provided is a method for using a stabilization stage to drive transistors in first and second charge pump stages that are interconnected by an intermediate circuit node.

As is well known, charge pump circuits are widely used in many electronic devices to obtain voltage values higher than a supply reference by using capacitors as charge storage elements. A classic charge pump circuit normally comprises four capacitors suitably connected to four enable signals, or phases, of the charge pump circuit.
A four-phase charge pump circuit is able to operate in a range of frequencies equal to 10-50 MHz. However, this range is insufficient for some applications.
Charge pump latch circuits are known to be able to output signals at rather high frequencies (e.g., 100 MHz and more), with respect to the traditional four-phase charge pump circuits which reach their maximum at 50 Mhz. This is possible since the charge pump latch circuits operate with only two phases and only use low voltage transistors.
A conventional charge pump circuit of the latch type is shown in FIG. 1. The charge pump latch circuit 1 has an input terminal IN connected to a first voltage reference (in particular the supply voltage Vdd), and an output terminal OUT connected to a load, which is illustrated by the parallel of a load capacitor CL and a load current generator IL inserted between this output terminal OUT and a second voltage reference (in particular the ground GND).
The charge pump latch circuit 1 comprises a plurality of N charge pump stages CB1-CBN inserted in cascade between the input IN and output OUT terminals and connected to a first enable terminal FX and to a second enable terminal FN supplying respective first and second enable signals, or phases (which for the sake of simplicity likewise are indicated with the references FX and FN). In particular, the enable signals, or phases, FX and FN are complementary to each other.
The charge pump latch circuit 1 is thus a two-phase circuit, able to operate at rather high frequencies (e.g., 100 MHz and more), due to the presence of suitable latch circuits L1-LN inserted in each charge pump stage CB1-CBN.
In more detail, each charge pump stage CBN comprises at least one first pump capacitor CUp and one second pump capacitor CDown inserted between the first enable terminal FX and a first inner circuit node Up and between the second enable terminal FN and a second inner circuit node Down, respectively.
Each charge pump stage CBN also comprises a respective latch circuit LN that is inserted between the first inner circuit node Up and the second inner circuit node Down and that comprises Low Voltage (or LV) transistors. In particular, each latch circuit LN comprises at least one pair of N-channel MOS transistors and a pair of P-channel MOS transistors, suitably connected to form opposite pairs of CMOS switches having control terminals or gates connected to the enable terminals FX and FN, for the alternating connection (and the relative charge sharing) of the pump capacitors CUp and CDown.
Since the charge pump stages CB1-CBN are placed in parallel, in the operation of the latch circuit 1 consecutive stages work in phase opposition with each other and the enable signals, or phases, FX and FN are at any time complementary to each other.
It is interesting to note that the voltage value VInt on an intermediate circuit node INT between two consecutive charge pump stages is kept practically constant during the whole clock period. Moreover, the widening of the range of operation frequencies allows a reduction in the sizes of the charge pump capacitors CUp and CDown and thus of the area occupied by the latch circuit 1 as a whole.
However, while the use of these charge pump capacitors CUp and CDown for the turn-on and the turn-off of the transistors of the latch circuit 1 allows there to be obtained a remarkable circuit simplification for the latch circuit itself, on the other hand they constitute a limit in the supply of current to the circuit, in particular when the supply voltage Vdd becomes close to the threshold voltage of the transistors in the latch circuit 1 or when the same is requested to supply a greater and greater current value.
In general, charge pump latch circuits do not operate very well at very low supply voltages. It is in fact possible to verify that, when the supply voltage Vdd decreases, a charge pump latch circuit of the type shown in FIG. 1 does not succeed in supplying a sufficient load current IL.
To overcome this limit, a charge pump latch circuit has been provided with a stabilization circuit connected to the enable terminals FX and FN in order to ensure a correct biasing of the transistors during the charge sharing periods.
For example, one solution described in European Patent Application No. 06 425497.2 of the same Assignee proposes the use of supplementary capacities only when the supply voltage becomes close to the threshold voltage of the transistors or when the charge pump latch circuit is requested to supply a greater and greater current. These supplementary capacities, by biasing the gate terminals of the pass transistors, allow the correct turn-on and turn-off of the same up to the end of the charge sharing.
This solution, even though it is able to operate with supply voltages lower than 1V, loses the typical characteristic of the charge pump latch circuits, i.e., the capacity of turning on and off the pass transistors slowly so as to avoid stressing the thin oxide layers whereon the gate areas of such transistors are realized.
Moreover, during switches at high frequency some asymmetries occur in the biasing voltages of the pass transistors between the upper and the lower circuit portion of the latch circuit.
These two drawbacks could cause stress on the oxides and an improper operation of the charge pump latch circuit, so as not to allow the meeting of the current and voltage requirements.
FIG. 2 shows a view of the charge pump latch circuit described in European Patent Application No. 06 425497.2. In this circuit, the phases FX and FN are, at any time, complementary to each other.
There are two different equivalent conditions as follows.
1) FN=Vdd, FX=0, VcD(i-1) is brought high, VcU(i-1) is brought low, VcD(i) is low while VcUpi is high. In this case the node PgU is brought high (at Vi) while the node PgD is bought low (by a Vdd with respect to a Vi), thus allowing the turning on of the transistor MpD and the turning off the transistor MpU. In the same semi-period the node NgD is brought high (by a Vdd with respect to Vi) and the node NgU is brought low. Thus the transistor MnD is turned on, while the transistor MnU is turned off. In this way, the charge pump latch circuit of FIG. 2 allows there to be obtained the charge sharing between the two adjacent bottom capacities, connected to the nodes VcD(i−1) and VcD(i), while it is forbidden between the two upper capacities, connected to the nodes VcU(i−1) and Vcu(i).
2) FN=0, FX=Vdd. In this case the latch circuit of FIG. 2 operates in a similar way to the previous situation. However, this time the charge sharing occurs between the two upper capacities, connected to the nodes VcU(i−1) and VcU(i).
It is to be noted that during the two semi-periods in which the nodes PgD, PgU, NgD, NgU, go high or low, they always remain constant for the whole time of the sharing, and at such a voltage as to allow an efficient turn-on and turn-off of the pass transistors connected thereto, also when the supply voltage Vdd is equal to 1V, as shown in FIG. 3. In particular, FIG. 3 shows the semi-period in which the charge sharing occurs between the two bottom capacities, i.e., the capacities connected to the nodes VcD(i−1) and VcD(i).
The transistors Mp1, Mp2, Mn1, and Mn2 allow the capacities Cb to update at the voltage they refer to and to recover that charge lost to drive the pass transistors.
To completely evaluate the characteristics of this circuit, we will evaluate the differences of potential between the various nodes. First, consider the voltage differences between the nodes of a generic stage j of the single charge pump latch circuit, as shown in FIG. 4 or equivalently of one of the charge pump stages of the circuit shown in FIG. 2.
Starting from the most intuitive condition, which is that of end sharing ("End Sharing" instant—FIG. 4), when the voltage difference between the nodes A and B is obviously void:
      V          AB              End                _                Sh              =  0During the sharing, the voltage at the node A has been reduced by ΔV, being:
      Δ                          V    =                    I        L                  T        C  while that at the node B increases by the same amount, thus at the beginning of the sharing there occurs:
      V          AB              Start                _                Sh              =      2        Δ                          V
But going back to the end sharing instant, at this point the phases switch ("Start No Sharing" instant) and, in consequence, the voltage at the node A decreases by a value equal to the supply voltage Vdd, while that at the node B increases by the same amount. This implies:
      V          AB              Start                _                NoSh              =      2        Vdd
During this semi-period there is no sharing in the stage j, while it is present in the adjacent stages, thus the capacity at the node A receives charge from that of the previous stage and the capacity at the node B transfer charge to that of the successive stage. At the end of the semi-period ("End No Sharing" instant) the voltage at the node A has increased by ΔV while that at the node B has been reduced by ΔV, thus there occurs:
      V          AB              Endt                _                NoSh              =            2            Vdd        -          2            Δ                                      V
The results so far obtained are summarized in the following table.
TABLE 1      V          AB              End                _                Sh              =  0End sharing between A and B       V          AB              Start                _                NoSh              =      2        Vdd  Start no sharing between A and B       V          AB              Endt                _                NoSh              =            2            Vdd        -          2            ΔV      End no sharing between A and B       V          AB              Start                _                Sh              =      2        ΔV  Start sharing between A and B
Now we will reconsider a double charge pump circuit, making reference to the schemes of FIGS. 5 and 6. In particular, the circuit comprises two upper nodes A and B and two lower nodes C and D whereto respective capacities are connected. It is verified that the potential differences between the nodes during the operation are the following ones.
1) under the Start Sharing Up and Start No Sharing Down conditions of FIG. 5:VBA=2ΔV; VCA=Vdd+ΔV; VCD=2Vdd; VDB=−(Vdd+ΔV)
2) under the End Sharing Up and End No Sharing Down conditions of FIG. 6:VBA=0; VCA=Vdd−ΔV; VCD=2Vdd−2ΔV; VDB=−(Vdd−ΔV)
Considering then the circuit complete with the switches as shown in FIG. 2, the absolute values of the voltages at the nodes at the end of the semi-period in which the charge sharing has occurred between the two bottom capacities are summarized in the following table.
TABLE 2VCUi-1Vi − Vdd + ΔVVCDi-1ViVCUiVi + Vdd − ΔVVCDiViPgDVi − VddPgUViNgDVi + VddNgUVi
Therefore, observing the differences of potential, the two pass transistors MpD and MnD on the bottom of FIG. 2 are on while those MpU and MnU on the top are off. Given the difficulty of obtaining perfectly crossed phases of the charge pump circuit, they are preferably realized overlapped or disoverlapped. In fact, when the phases are overlapped (i.e., when both FX and FN are equal to Vdd), the values of the voltages at the nodes become the following.
TABLE 3VCUi-1Vi + ΔVVCDi-1ViVCUiVi + Vdd − ΔVVCDiVi + VddPgDViPgUViNgDVi + VddNgUVi + Vdd
It results that the pass transistor MpU is on, but also that the transistors Mn1 and Mn2 are on. In consequence, the capacities Cb which drive the N-channel MOS transistors are unloaded on the pumping capacity C and in a different way with respect to each other. This creates an asymmetry of the potentials at the nodes between the high part of the charge pump circuit and the low part during the charge sharing phases, as graphically shown in FIG. 7.

Electronic circuit for generating low voltage and high frequency phases in a charge pump


A method for treating liquid waste water sludge after its generation and separation from the waste liquid involves forcibly impacting sludge particles and microbial cells against a contact surface to thereby fragment the particles into smaller-sized particles and break apart the organic cells. As a result of such a fragmentation process, water located in the pores of the sludge particles is released, thereby enhancing a subsequent dewatering process performed on the sludge. The destruction of the cells permits a subsequent solids digestion process to work more efficiently, while also permitting a more complete conversion to dissolved sugars, proteins and carbon dioxide.

Commonly, waste water is treated to remove contaminants that would otherwise prevent reuse of the water in potable water systems or that would contaminate receiving waters. Similarly, other waste liquids are treated to remove contaminants that might prevent reuse of the liquids in industrial processes. Such contaminants can include solid particulate matter, or dissolved acids, metallic salts and organic alcohols, as well as dissolved odorous gases. In the case of solid particulate contaminants, the contaminants are separated out of the liquid as a sludge.
Primary treatment of waste water often involves discharging the waste water into a primary settling tank. The water containing suspended solids is then allowed to remain in the settling tank for a period of time (detention time) sufficient to permit the solid particulates to gravitationally settle to the bottom of the tank as a sludge. This sludge is periodically or continuously removed from the bottom of the tank.
In another type of treatment (i.e., air flotation thickening), air bubbles are injected into the waste water prior to its admission into a subsequent settling tank so that the air bubbles can attach to the solid particulates to provide a lighter particulate condition. In the settling tank, the lighter particulates collect on the liquid surface, where they are skimmed off by an overhead travelling belt-type collector as a thickened sludge. The heavier particulates settle to the bottom of the tank as a second separate sludge.
Sludges can also be produced in a secondary treatment process as a by-product of an aerobic aeration process, wherein an aeration unit is included in aeration basins to bring gaseous air (oxygen) into contact with the particulates. Waste water flows into the aeration basins where it remains for a specified detention time. The air and incoming microorganisms are consumed by microorganisms in the mixed liquor in the aeration basins to produce a microbial floc, i.e., activated sludge. The mixed liquor from the aeration basins then flows to secondary settling tanks. The floc settles to the bottom of the secondary settling tanks, where it is periodically or continuously removed from the effluent liquid phase.
Another process involves anaerobic digestion, wherein the waste water is placed in a closed tank. The closed system allows anaerobic biochemical processes to take place, whereby some of the acids in the waste water are converted to methane gas and carbon dioxide, and solid particulates are thus consumed, or changed, to a liquid and gaseous form. A quantity of sludge remains as a by-product of the process. Typically, the tank is heated to an elevated temperature, at or near 95.degree. F., in order to keep the process ongoing. The methane gas is removed through a vent in the top of the tank, and the sludge is removed through a valved port in the tank wall.
Sludge generated as a result of any of the aforementioned waste water treatments is commonly disposed of by incineration, landfill operations, or ocean deposition. In some cases, the sludge is converted into fertilizer for agricultural use. A common problem associated with economical disposal of the various waste water sludges is the high water content of the sludge removed from the settling tanks. Typically, the sludge will have a water content of from about ninety percent to about ninety five percent, with only five to ten percent of the sludge mass being solids. This fact is an obstacle to the economical disposal of the sludge, since it increases the total sludge quantity that has to be handled. The high water content also makes the sludge somewhat more difficult to handle, because of its larger volume.
In an effort to increase the solids content in the sludge, the raw separated sludge has been subjected to various dewatering treatments, including vacuum filtration, centrifugations, pressure filtration, and drying on sand beds. However, these post-formation dewatering treatments have only been partially successful. The sludges resulting from such dewatering treatments still have only relatively small solids contents Typically, the solids content of so-called dewatered sludge is only about twenty to forty percent, with the remaining eighty percent to sixty percent being water.
The high water content of so-called dewatered sludge is due to the fact that much of the water is contained within the pores of the solid particles, or within narrow channels formed between agglomerated particles. Each particle is a relatively porous cellular structure that acts somewhat like a sponge to absorb and retain water molecules within the particle mass. Conventional dewatering processes are relatively ineffective in drawing, or extracting, water from the pores formed within such solid particles.

Method for treatment of waste water sludge


The present invention relates to a dishwasher which comprises: a washer main body in which a washing basin is provided to form a space in which dishes are washed; multiple fixed jet nozzles which are fixed to the inner wall of the washing basin and jet water for washing toward the dishes at a fixed position; and a horizontally arranged cylindrical rotating rack rotatably disposed in the washing bath, wherein the inner space thereof is split into multiple dish mounting portions where the dishes can be mounted and the rotating rack rotates along the periphery of the multiple fixed jet nozzles.

A typical dishwasher as described in Korean Patent Publication Gazette Number 2002-0077643 is characterized in that in terms of its structure and operation a fixed type dish-mounted rack onto which dishes (including containers or bowls) are mounted for washing is arranged in a washing basin, and the mounted dishes are washed while multiple rotational jet nozzles configured to spray water rotate around the fixed type dish-mounted rack.
The above-mentioned fixed type dish-mounted rack may be a fixed type quadrangle rack or a rack mounted on a turntable which rotates leftward or rightward.
As described above, the prior dishwasher wherein the fixed type dish-mounted rack is fixed and dishes are washed while rotational jet nozzles rotate, has the following problems.
First, a space applicability is bad because the arrangement of dishes is limited due to the structural limitation in the fixed type dish-mounted rack, as a result of which a bulky space is required for washing.
Second, it needs to mount the dishes onto the fixed type dish-mounted rack in a state that a concave side of each of the dishes is positioned to face downward in order to let washing or rinsing water to trickle down well, which makes the use of the dishwasher complicated.
Third, it may take relatively longer time to wash dishes because of the structure that the dishes are washed while rotational jet nozzles rotate around the fixed type dish-mounted rack.
Fourth, there may be a difference in washing force varying based on the positions or shapes of the mounted dishes, thus partially leaving unwashed food residues behind even after washing.
Fifth, it needs to additionally install an opening and closing door at a front side or a lateral side due to its structure, so mounting or unmounting dishes causes inconvenience when in use along with complicated mechanical structures.
Sixth, mold may grow due to incomplete dry in the dishwasher, so hygiene problems occur.
As listed above, in the prior dishwasher, a variety of problems occur due to the structural limitations, more specifically, due to the limited structure that a fixed type dish-mounted rack is fixed in the washing basin, and the rotational jet nozzles are configured to rotate. In order to improve the above-listed problems, it is urgent to develop a new dishwasher.


The present invention provides a liquid seasoning, which brings on a sufficient salty taste although it has a low common salt concentration, and which exhibits for example a pharmacological effect such as an antihypertensive effect at a high level. The present invention relates to a liquid seasoning, which contains the following components (A) to (C):

1. Field of the Invention
The present invention relates to a liquid seasoning.
2. Background of the Invention
Nowadays, interest in physiological functions of various components contained in food products is highly increasing. The Ministry of Health, Labor, and Welfare in Japan has certified food products containing components associated with such physiological functions or biological activities as Food for Specified Health Use (FOSHU). Such food products have been commercialized in the form of a beverage, yogurt, soup, miso soup, a prepared food product such as hamburger steak, a tablet-form confectionery, a tablet, or the like. It has been recommended that such food products be ingested once or twice a day.
Various materials have been proposed as materials having physiologically active functions. An example is a food material having an antihypertensive effect. Among others, peptide, γ-aminobutyric acid, chlorogenic acid, or the like are considered to be substances that are contained in food products and are highly safe. Thus, techniques of increasing the contents of such substances in food products or adding such substances to food products have been proposed (JP-A-06-78716, JP-A-06-197727, JP-A-2000-60489, and JP-A-2004-147560).
Since ingestion of common salts has adverse effects on renal diseases, cardiac diseases, and hypertension, addition of these to food products containing a high content of salts has been proposed. There have been a large number of techniques regarding the combined use of such a food material having an antihypertensive effect with soy sauce as a representative example of the aforementioned food products (JP-A-2002-87977, JP-A-2004-290129, JP-A-2004-194515, JP-A-2004-290088, JP-A-06-133723, and JP-A-02-167052). In order to ingest an effective amount of such a food material, a large amount of the food product should be ingested. However, ingestion of a large amount of such a food product results in ingestion of large quantities of salts. This leads to a decrease in the effect of ingesting the aforementioned food material, and thus it is unfavorable. Further, addition of a large amount of the aforementioned food material may affect the flavor of the food product.

This patent is subject to a terminal disclaimer.

Low salt liquid seasoning with antihypertensive activity

Abstract: (A) 3.55% or less by weight of sodium;

Abstract: (B) 0.5% to 4.2% by weight of potassium; and

Abstract: (C) 0.1% to 10% by weight of a food material having an antihypertensive effect,

Abstract: wherein the remaining portion of the liquid seasoning excluding component (C) is adapted by incorporation of a substance other than component (C) to have a nitrogen content (D) of 1.6% or more based on the total weight of the liquid seasoning.


Apparatus for testing the integrity of an access chamber system during or after installation thereof, includes a connection for making substantially fluid-tight connection with any part of the chamber system, evacuation device for evacuating the air within the chamber system to provide a vacuum therein, and an indicator for monitoring the level of vacuum over a predetermined period of time.

The present invention relates to apparatus and method primarily intended for testing the integrity of an access chamber system, at any stage prior to the completion of the installation, whereby there is an opportunity to correct any faults. The apparatus and method of the present invention may also be used on an existing access chamber system installation, to ensure that the installation is functioning correctly and/or to identify any faults or weaknesses in any part of the installation. This invention is designed to assist environmental issues of containing petroleum spirit (and associated vapor) within a chamber system and at the same time achieving a water free environment within said chamber.
The phrases "chamber system" or "access chamber system" are each to be interpreted as meaning any or all of access chamber itself, corbel, frame, neck or any extension piece(s), lid, suction lines or other utilities.
The installation of an access chamber involves a number of stages. Firstly, the chamber is fixed and sealed to the storage tank and this may be by one of the following:
(a) to the neck or lid of the manway lid, PA1 (b) to a separate upstand, or PA1 (c) to the fill pipe and/or other pipe flanges.
Then, one or more connections have to be made through the side wall of the access chamber for supply lines, electrical cables, and other utilities, and it may also be necessary to add an extension to the chamber, all of these connections must likewise be sealed. Finally, a corbel unit may be secured and sealed to the upper region of the access chamber, and this is also secured to a frame which is seated flush with ground level (the term "access chamber installation" is to be interpreted as encompassing any or all of these components).
At any or all of these stages, faults can occur in that the respective seals are not properly made. In particular, the interior of the access chamber must be adequately sealed to contain any spillage of petroleum spirit and/or associated vapor as well as against possible ingress of water from the ground surrounding the access chamber--the pressure on the chamber due to this ground water can be considerable, particularly in winter or wet weather conditions when the water table rises.
It is therefore the aim of the present invention to provide testing apparatus and method of testing which seeks to identify problems with the various seals used in the installation, ideally prior to the various stages of backfilling and concreting when remedial action would involve re-excavation, the testing apparatus conveniently being portable and suitable for a range of both round and square chambers of varying sizes.

Testing of access chamber system installation


A robot and a control method thereof. The robot has plural robot arms, each having at least one joint unit and a hand, and the control method includes calculating in advance and storing a plurality of grasp poses to generate grasp routes, sensing respective orientations of an object, the robot arms, and an obstacle, selecting one grasp pose from among the plurality of grasp poses by judging a movable angle range of the at least one joint unit, whether or not the object collides with the obstacle, and self-collision between the robot hands based on the sensed respective orientations, and generating grasp routes using the selected grasp pose. In the control method, the plurality of feasible grasp poses is calculated in advance and then stored, thereby rapidly and stably performing grasping of the object.

1. Field
Embodiments relate to a robot, which naturally and stably grasps an object, and a control method thereof.
2. Description of the Related Art
In general, machinery which performs motions similar to those of a human using an electric or magnetic action is referred to as a robot. Robots at the early stage were industrial robots, such as manipulators or transfer robots, which performed dangerous work, simple repeated work, and work requiring great strength in place of humans for the purpose of work automation or unmanned work on industrial lines. Recently, research and development of humanoid robots which have an external appearance similar to that of humans and coexist in working and living spaces of humans to provide various services is progressing.
Such a humanoid robot is provided with robot hands to grasp an object so as to achieve interchange and cooperation with humans in daily life, and the robot hand includes a plurality of fingers and a palm similar to a human hand. The robot hand adjusts stiffness of tips of the fingers through impedance control, thus achieving both precise work and flexible and stable work.
In a conventional robot grasp control method, the Inverse kinematic equation is solved whenever grasping is performed, and if a solution to the Inverse kinematic equation is not obtained, it takes a long time to generate grasp routes. Further, an alignment criterion of grasp poses, calculated by solving the Inverse kinematic equation, is inadequate and thus unnatural grasping may be performed.


A sub image display control system is composed of two sub image display control portions that are a display region control portion and a display timing control portion and a display output portion. A sub image data packet is composed of three types of data that are a sub image header, sub image display data, and sub image control data. In this structure, display spatial loss and display temporal loss can be remarkably reduced and sub images with wide applicability can be accomplished.

1. Field of the Invention
The present invention relates to an image process apparatus for superimposing a sub image such as a text on a main image such as a TV image, a method for forming a sub image data packet, and a recording medium thereof.
2. Description of the Related Art
Conventionally, translated texts for movies, control and adjustment indications for TV remote control operations, and so forth are superimposed on main images by sub image process systems.
The sub image display systems can be roughly categorized as character code systems and bit map data systems.
In the character code type sub image process systems, as shown in FIG. 27, character data (for example, characters and patterns) stored in a character memory 80 are retrieved from a character generator 81 as character code assigned to characters. The character code is transmitted to a sub image display portion 83. The character code is temporarily stored in a sub image frame buffer 84 and then required characters are displayed. In this system, dedicated hardware such as the character generator 81 is required. Since data to be transmitted to a display portion is only character code, the amount of data to be transmitted and display time for display process can be reduced.
However, since only prepared characters are displayed, the applications for displaying sub images are limited. Thus, although this system is suitable for a sound volume control screen for a TV or the like, it is not suitable in other contexts, such as superimposed texts on movies.
On the other hand, in the case of the bit map data systems, since bit map data of sub images are directly transmitted to the display portion, the dedicated hardware that generates sub images corresponding to the character code is not required. In addition, since the shapes of sub images that can be displayed are not limited, the number of applications for displaying sub images is great.
However, the data rate of the bit map data systems is much higher than that of the character code system. Thus, it can be said that the bit map data systems have losses of data transmission time and display process time.
In other words, in the bit map data systems, each pixel should have sub image color data, sub image contour color data necessary for superimposing a sub image on a main image, and a mixing ratio of the main image and sub image. Thus, the data rate of data to be transmitted to the display portion becomes large.
In the bit map data systems, since pixel data with the same amount of data for one screen of a main image should be transmitted to the display portion, display spatial loss takes place. Hereinafter, the data for one screen is referred to as a frame.
In addition, in the bit map systems and the character code systems, even if the shape of a sub image that is being displayed does not change, sub image data should be continuously transmitted at each display frame interval. Thus, in data transmission and display control of these systems, there is temporal loss.
Consequently, in the above-described conventional sub image process systems, the bit map data system is superior to the character code system in effectiveness of generating sub images. However, the bit map data system has display spatial loss and temporal loss.

Method for forming sub image data packet including data of sub image superimposed on main image, recording medium for recording sub image data packet, and image process apparatus


An improved flowable fill material and method for producing same that provides a medium for disposal of waste wood or other reclaimed debris, including wood treated with anti-microbial and/or anti-fungal materials including those containing copper, chromium and arsenic bearing compounds. The flowable fill material, comprised entirely of recovered and/or recycled materials, provides a medium for the safe disposal of chipped wood or other recovered waste by effectively encapsulating wood chips and substantially retarding further decomposition. Encapsulation of treated wood substantially reduces leaching of arsenic and/or other toxic substances contained therein. The systematic mixing of water, Class C fly ash, and wood chips in specific proportions carries out the present invention.

The present disclosure relates to the field of waste recovery. In particular, the production of a flowable fill material that incorporates waste materials. More particularly, the disclosure relates to rapid-setting and self-hardening flowable fill material that utilizes both waste wood chips and coal combustion by-products and provides a disposal medium for waste treated wood.
As world populations continue to rise, the demand for food, fuel and other needed resources rises commensurately. Ever-increasing consumption and corresponding production continues to stress global ecosystems at unprecedented levels. These facts act to fortify the compelling notion of recycling and waste reduction. The disclosure contained herein is comprised of recycled and/or recovered materials.
Scrap Wood represents a significant portion of many waste streams. Studies have shown that wood comprises 20-30% of many waste streams entering municipal solid waste landfills and can exceed 30% at construction and demolition debris landfills. In 1998, the EPA estimated that 136 million tons of building-related debris was generated, with wood representing the largest single component. This fact becomes much more concerning when it is considered that a significant portion of the wood currently in use is treated, thus a significant portion of the wood waste stream is comprised of treated wood.
Wood is commonly treated with preservatives to inhibit fungal and microbial decay. Common preservatives include creosote and pentachlorophenol. But the predominant preservative over the past 10-20 years has been chromated copper arsenate (CCA) that introduces copper chromium and arsenic to the wood. Chromium and arsenic and toxic to mammals and copper is toxic to most aquatic life.
During the CCA treatment process, wood is first dried in a kiln. The dried wood is then placed in a pressure vessel where a vacuum is pulled and effectively applied to the wood's pore space. A CCA solution is then introduced into the vessel and pressurized causing the solution to enter the pore spaces. After the pressure is removed, the wood is the removed from the vessel and allowed to dry over a drip pad.
There are three types of CCA-treated wood: Type A, Type B, and Type C, with Type C being the most common. Type C CCA is composed (by weight) of 34.0% As2O5, 47.5% CrO3, and 18.5% CuO.
The American Wood Preservers Institute (AWPI) estimates that treating wood with CCA increases the usable life of wood by 7 to 12 times. The actual CCA content in treated wood is significant. It is estimated that the amount of arsenic associated with CCA-treated wood currently in service in the state of Florida alone is 26,800 tons. Though the desirable characteristics of treated wood are obvious, it has been observed that chemically treating lumber takes a perfectly useable, recyclable material from a renewable resource and renders it toxic.
CCA became widely used in the late 1970s and now represents approximately 80% of the treated wood market in the U.S. It is estimated that 6.5 billion board feet are treated each year in the U.S. Treated wood is particularly well suited for use in residential and small commercial decks, and estimates are that CCA treated wood has been used in over 80% of residential decks since the early 1970s. As the treated wood that is now in service expires in the coming years, the amount of CCA-treated wood arriving at landfills is expected to peak sometime between the years 2020 and 2030.
It is difficult to visually discern between treated and untreated wood, thus most states currently CCA-treated wood in the same manner as other discarded wood and wood products. Therefore, most treated wood in the waste stream ends up at unlined construction and demolition waste landfills. A far lesser amount is disposed in municipal solid waste landfills or incinerated at waste to energy facilities. Since the toxic metals in CCA easily survive the incineration process, burning treated wood is discouraged. Further, incinerating can convert trivalent chromium into highly toxic hexavalent chromium.
As previously discussed, the metals in CCA easily survive incineration and escape with flue gas. Further numerous studies have shown that CCA, particularly the arsenic component easily leaches from treated wood, thus, incineration, landfill disposal and processing into mulch are not suitable waste outlets for CCA-treated wood. Though little data is currently available, there is growing concern about the impact that CCA-treated wood in landfills may ultimately have on groundwater. Given these problems associated with CCA-treated wood and the existing and ever-growing problems associated with disposal or reuse, there exists a substantial need for an economical, safe, and environmentally responsible means of disposing of expired CCA-treated wood.
Concrete is one of the oldest and most important composite materials known to man and generally consists of a course aggregate (rock and/or gravel), sand, and hydrated Portland cement. The finished properties of concrete generally depend on several factors: ratio of cement, sand and aggregate; ratio of water to cement; nature of the course aggregate; mixing and laying methods; and curing time.
Portland cement, named after its likeness to the indigenous limestone of Portland Bill, England, is a mixture of primarily four minerals: tricalcium silicate (3CaO.SiO2); dicalcium silicate (2CaO.SiO2); tricalcium aluminate (3CaO.Al2O); and tetracalcium aluminoferrite (4CaO.Al2O3.Fe2O3). Typical composition, expressd in terms of oxides, is 65% Ca, 20% SiO2, 5% Al2O3, with the balance comprised of Fe2O3 and other admixtures. Portland cement is generally obtained by mixing materials which supply lime (usually limestone or chalk) with materials which supply silica and alumina (usually clay) and firing these mixtures to ˜1200 K for some period of time. Hydration of the oxides in Portland cement is a reaction that proceed for a very long time, as evidenced by concretes continued measurable hardening for years.
Fly ash is the inorganic non-combustible portion of coal that remains after pulverized coal is burned, and is generated in huge quantities by coal-fired electric generating facilities throughout the U.S. Fly ash is comprised of glassy, spherical shaped particles that are typically recovered from flue gas by means of electrostatic precipitators. There are two common Classes of fly ash: Class C and Class F. Class C is produced from burning lignite and sub-bituminous coal, and Class F is produced from burning anthracite and bituminous coal. The comparative typical mineral content of Class C fly ash, Class F fly ash, and Portland cement is shown below in Table 1.
TABLE IChemicalClass CClass FPortlandCompoundFly AshFly AshCementSiO39.9054.9022.60Al2O316.7025.804.30Fe2O35.806.902.40CaO24.308.7064.40MgO4.601.802.10SO33.300.602.30Na2O & K2O1.300.600.60
An important characteristic of Class C fly ash is self-cementing or self-hardening when mixed with water, characteristics enabled by a relatively high CaO (lime) content. As can be seen in Table 1, the CaO content of typical Class C fly ash is almost three times that of Class F.
Fly ash has found particular use as an additive in concrete and pavement materials. When added to Portland cement, fly ash effectively reduces the amount of lime needed and aids in converting free lime to calcium silicate hydrate, a substance similar to the dicalcium and tricalcium silicates—the strongest paste portion of concrete—formed during curing. Concrete enhanced with Class C fly ash will typically strengthen faster than plain concrete due to the additional lime content. Concrete enhanced with Class F fly ash will typically strengthen slower than plain concrete until about 7 days, then it will typically strengthen at a faster rate.
The environmental benefits of using fly ash are numerous. Materials such as Portland cement, lime and crushed stone require energy to produce. Utilizing one ton of fly ash to replace an equivalent ton of any of these materials conserves the energy equivalent of one barrel of crude oil, as well as reducing CO2 emissions by approximately one ton.
A Flowable fill prepared from fly ash, also commonly known as controlled density fill, has many highly desirable properties. It is a fill material that flows easily, is self-leveling, self-compacting, and non-settling after hardening in place. A fly ash-based flowable fill will easily encapsulate whatever has been placed in a trench or void and will provide protection after hardening. A fly ash-based flowable fill can typically be placed in any weather at any time and will displace standing water. Fly ash-based flowable fill can typically be prepared using conventional concrete mixing equipment and pumped using conventional concrete pumping equipment.
The disclosure contained herein takes advantage of the self-cementing properties of Class C fly ash and enhances these properties by the addition of brine to produce a flowable fill material which not only provides a medium for a preferred disposal means for treated wood, but also provides a highly effective fill material useful for a plurality of backfilling applications.

Flowable fill and flowable fill method for disposal of recovered waste


A new and distinct Clematis hybrid named ‘TRA35’ is disclosed, characterized by a very large flowers, up to 22 to 30 cm in diameter. Flowers have uniquely elontaged tepal shapes, and a violet-blue coloration, with a pink median streak. The new cultivar is a Clematis typically suited for ornamental container and landscape use.

The new cultivar is the result of a breeding program, with the objective to develop new varieties of Clematis with strong growth and interesting flowers for commercial purposes. The new variety was discovered as a single plant among a group of seedlings of unknown parentage observed by the inventor, Bernard Travers. The discovery was made during 2004. Discovery was made at a commercial nursery in Orleans, France.
After selecting and isolating the new cultivar, asexual reproduction of the new cultivar 'TRA35' was first performed at a commercial nursery in St. Jean Le Blanc by vegetative cuttings during 2005. The inventor maintained the original plant, and observed the clones for several years, to establish stability of the characteristics. 'TRA35' has since produced multiple generations and has shown that the unique features of this cultivar are stable and reproduced true to type. Public access to plant material first occurred during early Spring of 2014. Before this time, plant material was under the control of the inventor, or those directly authorized by the inventor. This variety is described in CPVO application no A201302537. A copy of this application and technical question are included with the mailing of this application. Page 5 of the application verifies at the time of the EU application sales had not occurred either in the EU or any other country. Neither the EU application, nor any other publication resulted in an enabling disclosure of 'TRA35'.