An apparatus for preparing blended beverages, particular frozen drinks, in which an ice dispenser, liquid mix dispenser and blender unit are combined in a single unit. A programmable timing mechanism regulates the period of time the ice dispenser, liquid mix dispenser and the blender unit operate, so that an operator needs to only push momentary push buttons to effect operation of the apparatus. The apparatus will automatically deliver an appropriate amount of ice and liquid to the blender unit and will turn on the blender at the appropriate time and for the appropriate amount of time to prepare a frozen drink of desired size.

The demand for "frozen drinks" has soared. These drinks, delicious concoctions having exotic names conjuring images of tropical islands, have become the refreshment of choice for a large portion of the adult population. Preparing these drinks, however, has never been an easy task. Typically, several ingredients including liquor must be added in a time consuming process to a finely-divided ice/water mixture having a thickness and almost creamy texture prior to serving. In the busy atmosphere of a bar, club, or restaurant setting, where profit is dictated largely by controlling portion size, waste connected with the preparation of such drinks has never been easy to control.
Typically, a frozen drink is prepared by placing the drink ingredients, including ice, into a conventional blender. The blender is then operated to reduce the ice to a pulverized state, and the resulting mixture then transferred to a glass for consumption by a patron. Because of the large span of time required in measuring and mixing the ingredients, several portions of a given drink are usually prepared in a blender simultaneously. If these drinks remain unsold at the time when the ice within the mixture has melted into a liquid then the drinks must be discarded.
Aside from the loss of product associated with conventional blending methods, an additional problem is frequently encountered with mixed drinks; the amount of time consumed while the blender grinds the ice limits the number of customers that a bartender can serve. Most bartenders cannot receive customer orders when their attention is focused upon a noisy blender. The reduced sales caused by conventional blending techniques can have a significant long-term economic impact upon the proprietor of an establishment serving mixed drinks.
In the past, machines for preparing frozen drinks have been developed. For instance, U.S. Pat. No. 4,681,030, issued Jul. 21, 1987 to John M. Herbert discloses an apparatus in which an ice-shaving machine is combined with a blender in a single unit. The ice shaver and blender are electrically wired together, and programmable timing means are provided. However, the apparatus provides no means for dispensing liquid mixes which must be manually added with any liquor to the blender. Additionally, U.S. Pat. No. 4,653,281, issued Mar. 31, 1987 to Richard F. Van Der Veer discloses a drink making method and apparatus for making beverages having dispensed ice-slush and a blend of additional flavorings. The system includes a blender for mixing the constituents, and the device dispenses ice into the blender while liquor and fruit flavorings are manually placed therein. Also, U.S. Pat. No. 3,156,103, issued Nov. 10, 1964 to Anthony J. Ross teaches an apparatus for dispensing a beverage into a cup and a measured charge of ice into the cup for cooling the same. The apparatus includes an ice producing and storing device, and a plurality of dispensers for more than one type of drink, with the dispensers electrically operative to dispense a premeasured charge of liquid.
Further, U.S. Pat. No. 4,993,593, issued Feb. 19, 1991 to Ralph Fabiano and Fortunato Colacino discloses a device for mixing and then dispensing hot liquids in a premeasured amount.
None of the above inventions and patents, taken either singly or in combination, is seen to describe the instant invention as claimed.


Methods for fabricating interconnect structures implementing air gaps therein is provided. In one embodiment, a semiconductor substrate with a first barrier layer formed thereon is provided. A first dielectric layer is formed above the barrier layer. The first dielectric layer is thereafter patterned and etched to form a plurality of stakes having first openings therebetween, the plurality of stakes for providing mechanical supporting strength for the interconnect structure. A sacrificial layer is formed in the first openings and above the plurality of stakes. A hard mask layer is formed above the sacrificial layer. A light sensitive layer is formed over the hard mask layer and is thereafter patterned to define a pattern therein. The hard mask layer, the sacrificial layer, and the first barrier layer are etched according to the pattern in the light sensitive layer to form second openings. The second openings are filled with a conductive material to form metal lines. The sacrificial layer is thereafter removed to form air gaps. A second barrier layer is formed over the plurality of stakes and the metal lines, the second barrier layer thereby sealing the air gaps.

The present invention relates generally to the fabrication of semiconductor devices, and more particularly to methods and structures that implement air gaps in semiconductor devices.
As device dimensions continue to shrink, a reduction in interconnect linewidths leads to increased line resistance (R) for signals. Further, reduced spacing between conducting lines creates more parasitic capacitance (C). The result is an increase in RC signal delay, which slows chip speed and lowers chip performance.
The line capacitance, C, is directly proportional to the dielectric constant, or k-value of a dielectric material. A low-k dielectric reduces the total interconnect capacitance of the chip, reduces the RC signal delay, and improves chip performance. Lowering the total capacitance also decreases power consumption. The use of a low-k dielectric material in conjunction with a low-resistance metal line provides an interconnect system with optimum performance for the ULSI technology. For this reason, prior art attempts to reduce the RC delays have focused on utilizing material with a low-k to fill the gaps between the metal lines. Silicon dioxide (SiO2) has been conventionally preferred as a dielectric material even though it has a relatively high dielectric constant (relative to vacuum) of about 4.1 to 4.5 because it is a thermally and chemically stable material and conventional oxide etching techniques are available for high-aspect-ratio contacts and via holes. However, as device dimensions decrease and the packing density increases, it is necessary to reduce the spacing between conductive lines to effectively wire up the integrated circuits. Therefore, a large number of lower dielectric constant materials are currently being investigated to reduce the RC value of the chip further. These include among many others fluorinated SiO2, aerogels, and polymers. Another method being proposed to lower the dielectric constant even further is to form air gaps between the interconnect lines. While silicon dioxide has a dielectric constant of about 4 and greater, the dielectric constant of air is about 1.
The formation of air gaps between conducting lines of high speed integrated circuits (IC's) is typically a combination of the deposition of a porous dielectric layer, selective etching of the dielectric layer, the deposition of a metal layer, and the removal of the dielectric layer to form the desired air gaps.
U.S. Pat. No. 5,407,860 (Stoltz et al.) discloses an air gap process by etching low-k material out between metal lines and forming a dielectric layer thereover.
U.S. Pat. No. 5,444,015 (Aitken et al.) shows a method for forming air gap between metal lines by removing the dielectric material between the lines.
U.S. Pat. No. 5,324,683 (Fitch et al.) shows a method for forming an air gap between metal lines by forming a dielectric layer between metal lines, forming an etch barrier layer(s) thereover and opening a hole in the etch barrier layer and isotropically etching the dielectric layer for form air gaps.
Although air is the best dielectric material for lowering the RC value, unfortunately the use of air gap structures in integrated circuit fabrication has been hindered with problems. Overall mechanical strength of the device is reduced correspondingly and lead to structural deformation and a weakened structure can have serious effect in various aspects of subsequent integrated circuit fabrication. To illustrate, FIG. 1 depicts a cross sectional view of a portion of an interconnect structure having air gaps formed therein according to a conventional method of interconnect formation. A semiconductor substrate 10 has a patterned barrier layer 20 formed thereon. A plurality of metal lines 80 are interposed between portions of the barrier layer 20 with air gaps 90 formed therebetween. A barrier layer 100 seals the upper surfaces of the plurality of metal lines 80 and air gaps 90. A dielectric layer 120 is subsequently formed above the barrier layer 100. While a portion of dielectric layer 120 overlying barrier layer 100 may not suffer from an overlay or dishing problem, portions of dielectric layer 120 overlay into a region 130 otherwise reserved for an air gap due to the insufficient support from underneath barrier layer 100. This overlay problem is particularly notorious in areas of the interconnect structure where metal lines are few and far between or where the distance between successive metal lines is wide. This overlay problem can cause many reliability problems. One such problem is distortion on the surface of the dielectric layer, which would require a planarization process to cure the distortion defect.
Accordingly, what is needed in the art is an air gap interconnect structure and method of manufacture thereof that addresses the above-discussed issues.

Air gap interconnect structure and method thereof

Taiwan Semiconductor Manufacturing Co., Ltd.


The invention relates to a forward-amplifying filter circuit which has an analogue filter with poles and zero points, having a first number of integrators which form the poles, which are arranged in series with one another and which each have an active component with a capacitance arranged in parallel therewith, having a second number of coefficients which form the zero points and which have at least one forward-amplifying path which contains a differentiating element, having a summing node whose input side is connected to the forward-amplifying paths of the zero points and whose output side is connected to the input of the last of the active components arranged in series. The invention also relates to an analog/digital converter circuit having such a filter circuit.

To produce various transfer functions, suitable filters are used. In this case, a distinction is drawn between active, i.e. analog, filters and digital filters. In the case of analog filter voltages, the processed signals are present in the form of continuous time functions. The functionality of an analog filter may also be provided in digital form. On account of the additional circuit complexity required for this and the fact that, particularly with a large number of data items to be processed, the analog/digital conversion and the digital/analog conversion are also time consuming, analog, continuous-time filters are used, particularly for very fast applications.
One special form of such an analog, continuous-time filter has a forward-amplifying architecture. Such forward-amplifying filter circuits, which are also frequently referred to as feedforward filters in technical literature, are known generally and are distinguished particularly by their performance. Forward-amplifying filter architectures are used particularly in telecommunications, for example in mobile radio and for broadband applications.
FIG. 1 uses a schematic block diagram to show the general design of a fifth-order forward-amplifying filter circuit. Such a filter generally comprises poles and zero points, the poles being produced by integrators a1–a5 and the zero points being produced by coefficients c1–c5. The coefficients c1–c5 are formed by forward-amplifying paths in FIG. 1. In addition, feedback paths with coefficients d1, d2 are provided. At the output of the forward-amplifying filter circuit, there is a summing node, which means that an analog output signal Vout can be provided at the output from the analog input signal Vin which is input on the input side.
FIG. 2 shows a known circuit implementation for the forward-amplifying filter circuit from FIG. 1. In this case, the integrators a1–a5 have been produced by operational amplifiers OP1–OP5, resistors R1–R5 and capacitances C1–C5, and the feedforward coefficients c1–c5 have been produced by resistors RK1–RK5. To produce the summing node at the output of the filter circuit, however, an additional summing amplifier OPout with a parallel-connected resistor Rout is required. However, an additional operational amplifier OPout not only increases the power consumption of the overall filter circuit but also means an increase in execution time. Although it would also be possible to dispense with the summing amplifier OPout at the output of the filter circuit, the filter circuit would no longer operate in isolation in this case.
As an alternative to the implementation in FIG. 2, the integrators a1–a5 would also be able to be formed by transconductance amplifiers (OTA, voltage/current amplifiers), and the feedforward coefficients would be able to be formed by transconductance amplifiers or capacitances. The transconductance amplifiers or capacitances which such a filter circuit requires in the forward-amplifying paths are synonymous with nonlinear forward amplification, however, which means that use is preferably made of filter circuits in which the integrators are produced by operational amplifiers, resistors and capacitances, as shown in FIG. 2.
In such power-optimized filter circuits, the number of poles should be as equal as possible to the number of amplifying components.


The urine collection insert, or infant urine collector, is an insert for placement inside a typical diaper. Placement of the insertion preferably occurs with the first diapering after birth. The infant's urine is channeled into a portal, through a one-way valve, and into a reservoir. Coincident with the next diaper change, the insert is collected and the urine withdrawn for testing. Throughout the process, all chain-of-custody requirements being followed. The design of the urine collection insert prevents contamination from the infant's feces, most likely to be meconium—the feces of a newborn formed from epithelial cells, lanugo (hair), mucus, amniotic fluid, bile, and water. By avoiding contamination, the integrity of the sample will be maintained and accuracy of results increased, whether positive or negative.

The opioid crisis has caused compilations that ripple beyond adult addiction. When a pregnant mother uses opioids, or other substances including alcohol or cocaine, the baby is born with his or her own addiction and therefore, anticipated withdrawal.
Treatment options exist for babies born with substance exposure issues, but early diagnosis is critical.
For example, if a newborn was exposed to heroin or cocaine, withdrawal will likely occur within forty-eight hours. Thus, treatment must begin quickly.
But if a newborn was exposed to methadone or benzodiazepines, withdrawal may be delayed as long as two weeks. Without early diagnosis, the infant has likely left the care and supervision of the hospital and been returned to a home where proper care may not be available. Thus, early diagnosis is critical to planning future care.
Current diagnostic methods require drawing blood from the baby. This is not only emotionally traumatic for the baby and parents, but the removal of blood can be taxing on a baby given its limited blood supply due to its small size.
Because the current testing methods require a blood draw, the diagnosis is only performed when there is suspicion that the mother has been abusing or using substances, or when the baby exhibits symptoms of withdrawal. This is an imperfect system, and thus some babies are not tested, and as a result treatment is delayed.
What is needed is a device for easily collecting fluid from a baby for testing purposes.


Session Inter-Device (SID) mobility networks (50, 100, 150) are described in which a seamless transfer of a communication session from a first device (56, 106, 116) to a second device (66, 116, 166) can be achieved without interrupting the active session. According to the SID mobility network (50), the transfer can be accomplished by transferring away from the Transferring Node or first device (56) the IP address associated with the active session (58) so that the network (50) will route the session to the desired Target Node or second device (66). The Transferring Node (56) transfers its IP address (58) to the Agent (60) and stops requesting data packets addressed to its IP address (58). The Agent (60) then begins to request and eventually receive the packets addressed to the Transferring Node's IP address (58). The Agent (60) then transfers the packets to the Target Node (66). In an alternate SID mobility network (100), the Transferring Node (106) transfers a session specific IP address (114) to the Agent (110). The Agent (110) then transfers packets sent to the session specific IP address (114) to the Target Node (120). In another SID mobility network (150), the Transferring Node (162) obtains a temporary IP address (170) and transfers its IP address (164) to a Session Agent (166). The Session Agent (166) begins to request and eventually receive the packets addressed to the Transferring Node's IP address (164), and for each received packet determines if it belongs to the session the Transferring Node (162) requested to transfer to the Target Node (176). If it does, the Session Agent (166) will transfer the packet to the Target Node (176) at the Target Node's IP address (178). If it does not, the Session Agent (166) will transfer the packet to the Transferring Node (162) at its temporary IP address (170). In each SID mobility network, the session with respect to the Correspondent Node continues without interruption throughout the transfer, thereby providing a seamless transfer of the session from a first device to a second device.

This invention relates generally to a method and apparatus for transferring a communication session from a first destination device to a second destination device in a computer network with minimal disruption to the communication session, and more particularly to a method and apparatus for transferring a communication session between a Transferring Node and a Correspondent Node in an Internet Protocol (IP) based network from the Transferring Node to a Target Node without disrupting the communication session with respect to the Correspondent Node.
Over the past decade, the number and size of computer networks, such as the Internet, Intranets, Wide Area Networks (WAN's), and Local Area Networks (LAN's), have increased dramatically. In addition, the number of users who access these networks from various locations over the network has grown. For example, many users are accessing their networks from other office buildings on and off corporate campuses, from their homes, from client offices, etc. The increase in the number and size of networks, as well as the increase in the number of mobile users has introduced several problems including the inability to provide uninterrupted access to the network for Mobile Nodes. One of the main problems confronted has been the inability to allow a user to disconnect from one IP sub-network in a large IP network and reattach to the large IP network from another sub-network without interrupting ongoing network communication sessions. For example, users may be required to manually change IP addresses and reinitialize sessions when they move about the network in order to be recognized as an authorized network user or to access their network files. A solution to this problem was presented in the form of an extension to the Internet Protocol called Mobile IP.
Mobile IP is a recently standardized method for enabling what this description terms Session Inter-Subnet (SIS) mobility in IP networks. SIS mobility refers to the ability of a device to disconnect from an IP network such as the Internet and reattach to the IP network on another sub-network without interrupting ongoing sessions. Mobile IP is documented in Internet Engineering Task Force (IETF) Request For Comment (RFC) 2002, 2005, and 2006. Textbooks such as "Mobile IP: The Internet Unplugged" ISBN 0-13-856246-6 by James D. Solomon also document Mobile IP. Mobile IP is one of the dominant IETF standards in mobility and operates as a highly secure and dynamic packet data device mobility service. It provides a network layer solution for moving active sessions with a device as the device moves between IP sub-networks. Therefore, the user does not have to manually change an IP address and reinitialize sessions when moving around within the IP network, nor do Correspondent Nodes need to take action to maintain communication with Mobile Nodes.
According to Mobile IP, a network connection can be broken down into a Correspondent Node and a Mobile Node connected via an IP network. The Correspondent Node is a node located elsewhere on the IP network that is currently communicating with another node called the Mobile Node. When the Mobile Node is at its home location, in other words directly attached to its home sub-network, it receives communications from the Correspondent Node via routers that direct the communications to a home router, which in turn directs this information to the Mobile Node. During this active communication session, the Mobile Node may wish to move to a different sub-network, called a foreign network, while continuing to maintain the active session throughout the move. In order to do so without disrupting the Correspondent Node and the current communication session, the Mobile Node must continue to be reachable via the IP address it has been using throughout its communication session with the Correspondent Node. This task is accomplished in Mobile IP by use of a Home Agent.
The Home Agent (HA) intercepts communications meant for the Mobile Node and forwards these communications to the Mobile Node wherever it currently is located. The HA accomplishes this by use of a low-level protocol within the Internet Protocol (IP) suite called Address Resolution Protocol (ARP). ARP associates or "binds" the physical address of a node on the local sub-network with its IP address. Multiple IP addresses can be associated in this fashion with a single physical address. By using ARP, the HA associates its physical address with the IP address(es) of the Mobile Node(s). Following this association, the HA will receive traffic (or data) sent from the Correspondent Nodes to the Mobile Node's home address and the Mobile Node will register a "care of" address with the HA identifying its current location. The Mobile Node keeps the HA aware of its current care-of address by sending registration messages to the HA. The HA then forwards communications to the "care-of" address on the foreign network via an IP tunnel. Various types of IP tunnels exist in the art (IETF RFC's 1701, 2003, 2004), but each is essentially used to encapsulate the original communication as an inner IP packet (inner packet) or payload section of another larger IP packet (outer packet) which is addressed from the tunnel starting point to the tunnel endpoint. The tunnel endpoint node removes the outer packet header information, restoring the original packet (or communication).
One possible endpoint for the tunnel carrying forwarded packets for the Mobile Node from the Home Agent is a Foreign Agent (FA). The FA is a router with at least one interface on the Mobile Node's current foreign sub-network. Upon receipt of the tunneled packet, the Foreign Agent can provide a care-of address for the Mobile Node so the packet will be properly delivered. The FA may also de-tunnel packets forwarded by the Home Agent for the Mobile Node by removing the outer packet header, thus restoring the packets back to the packets originally intercepted by the Home Agent. The FA then sends these de-tunneled packets directly to the Mobile Node on its current foreign sub-network via the Mobile Node's link-layer address.
Another possible endpoint for the tunnel from the Home Agent is the Mobile Node itself. In this case, the Mobile node may acquire a temporary IP address on the foreign sub-network. This temporary address, called a co-located care-of address, functions as the care-of address for the Mobile Node while connected to the foreign sub-network. The Mobile Node may acquire this address by any means for generating IP addresses known in the art such as a Dynamic Host Configuration Protocol (DHCP) which is documented in IETF RFC 2131. Using this method of acquiring an address, the Mobile Node contacts a DHCP server requesting an unused IP address on the Mobile Node's current foreign sub-network. DHCP is increasingly common on current IPv4 networks, and future networks employing IPv6 will allow for nodes to dynamically create IP addresses without needing to acquire the addresses from a server. Thus, co-located care-of address generation is an increasingly viable option for receiving tunneled packets from the Home Agent, freeing the Mobile Node from the constraint of only visiting foreign sub-networks that have FA's.
The operation of Mobile IP outlined above was designed to enable a specific type of mobility, but there remain other types of mobility that have not been addressed. Mobile IP partially fills the mobility gap in the IP protocol suite by accommodating device movement between IP network attachment points or sub-networks while maintaining all of the device's IP communication sessions. For example, Mobile IP enables a notebook PC to move from its office Ethernet connection to a wide-area wireless packet data network with all of its network communications sessions with multiple Correspondent Nodes in tact. This Session Inter-Subnet (SIS) mobility has become increasingly common, thereby increasing the need for Mobile IP. However, Mobile IP does not address the problem of moving active communication sessions between different devices. For example, a user may wish to transfer a video conference from a fixed desktop machine to a notebook PC or Personal Digital Assistant (PDA) as the user leaves the office in order to continue participating. This description terms moving active IP networking sessions between devices such as in the above example "Session Inter-Device" (SID) mobility.
Networks other than IP networks provide transfer features and can enable movement of active communication sessions between devices. However, these networks do not address the variety of media and devices present in IP networks. For example, the International Telecommunications Union (ITU) H.323 protocol suite provides for call transfer based on the traditional telephony model and can accommodate transfer of multimedia sessions, and the IETF Session Initiation Protocol (SIP) family (IETF RFC 2543) with proposed extensions (IETF Draft "SIP Call Control Transfer" draft-ietf-sip-cc-transfer-02) provides for call transfer accommodating multimedia sessions as well. However, both of these approaches operate at a level often requiring modification to applications, including applications running on Correspondent Nodes. Furthermore, traditional telephony call transfer methods involve a gap or "hold" in the communication session while the call is redirected.
Thus, there is a need for a system and method capable of providing selective transfer of active communication sessions between nodes on an IP network from one end device to another end device with minimal disruption of the selected and unselected communication sessions.

This patent is subject to a terminal disclaimer.

Method and apparatus for transferring a communication session


Systems directed to the art of setting a hook in a fish mouth. The systems have a sleeve substantially housing a plunger assembly and a biasing member, and an activation member assembly with a rotatable finger. The biasing member applies force between the plunger assembly engages and the activation member assembly when the system is engaged. A fish hook is attached via fishing lines to both the activation member assembly and the plunger assembly. Pulling the hook releases the plunger assembly from the activation member assembly and pulls the line attached to the plunger assembly, setting the hook in the fish's mouth.

Fishing is an activity predicated on patience and timing. When a fish snatches the bait offered on an unassuming hook or takes a bite of an attractive lure, the angler must be ready to set the fishing hook in the mouth of the fish to ensure the capture of the water faring beast, or else embrace the fate of being the losing party in another epic fishing saga. Setting a hook is generally accomplished by pulling back on the fishing rod in a quick jerking motion. If poorly timed, either too early or too late, the fish may escape as the hook has either not yet entered the fish's mouth or the fish has already taken the bait and swam away. Accordingly, anglers would benefit from a device capable of more consistently setting a hook in a fish's mouth at a more advantageous time.


In a circuit simulation system, a storage section is configured to store a circuit data, an analysis condition data and an output data. An initial data setting section reads out the circuit data and the analysis condition data from the storage section and sets an initial data and a convergence condition for a solution calculating process based on the circuit data and the analysis condition data. A processing section generates a circuit equation to each of a voltage variable and a current variable based on the circuit data, and executes the solution calculating process based on the initial data to calculate a solution. A convergence determining section executes a convergence determining process of whether or not the solution meets the convergence condition, on the voltage variable. An output section stores the solution into the output data when it is determined to meet the convergence condition. A repetition control section controls the processing section to calculate a next solution by using the solution as the initial data, when it is determined not to meet the convergence condition.

1. Field of the Invention
The present invention relates to circuit simulation of a semiconductor integrated circuit, and more specifically, to a circuit simulation method in a designing process of the semiconductor integrated circuit, and a circuit simulation system.
2. Description of Related Art
In recent years, the scale of a semiconductor integrated circuit has been made large and a time necessary to analyze a circuit is increasing in the design and development of the semiconductor integrated circuit. A circuit simulation program such as SPICE is frequently used for the analysis of the circuit. An analyzing method in the SPICE is well known, and is shown in the Web Site (http://embedded.eecs.berkeley.edu/pubs/downloads/spice/index.htm). Although there is no flowchart, a source code of the SPICE has been opened.
A circuit simulation method of analyzing a circuit at high speed by using the SPICE is disclosed in Japanese Patent Application Publications (JP-A-Heisei 10-222555, JP-A-Heisei 11-110420, JP-P2006-133994A. In Japanese Patent Application Publication (JP-A-Heisei 10-222555), a circuit is divided into a plurality of sectional circuits and analyzed, and in Japanese Patent Application Publication (JP-A-Heisei 11-110420), a circuit is analyzed in advance while a precision parameter is varied. Japanese Patent Application Publication (JP-P2006-133994A) discloses that a current value or a voltage value is calculated.
In the circuit simulation program represented by SPICE, a circuit equation is generated based on Kirchhoff's voltage law and current law to express connections in the circuit. The circuit equation is represented as a matrix, and the circuit is analyzed through calculation of the matrix. That is, based on data of elements contained in the circuit and data of connections of them, the equation F(X)=0 is defined, in which a solution X (=unknown variable) is a value of a voltage (voltage variable) at each of nodes in the circuit and a value of a current (=current variable) flowing through each of the elements.
An iterative solution method such as the Newton iteration method is generally used for calculation of the solution of the circuit equation. The Newton iteration method is a technique of calculating an approximate solution of a true solution satisfying the equation F(X)=0, by specifying a variable i indicating the number of times of the iteration, and repeating calculation of a value of the ith unknown variable X(i) by using a value of the (i−1)th unknown variable X(i−1). It is called convergence that an amount of change of the unknown variable X settles in a tolerance specified in advance, and a process is completed assuming that the solution has been calculated at that time.
FIG. 4 is a block diagram showing the configuration of a conventional circuit simulation system. Referring to FIG. 4, the conventional circuit simulation system includes a storage unit 1, an analysis processing section 2, and a calculation result outputting unit 7.
The storage unit 1 is exemplified by a hard disk drive and is connected to an I/O interface 4. The storage unit 1 stores a circuit data 10 specifying the circuit configuration such as data indicating elements contained in a circuit and connections of them, and an analysis condition data 20 as a setting data used for performing analysis such as an analysis start time, and an analysis end time. Moreover, an analysis result is stored in the storage unit 1 as a waveform data 40.
The analysis processing section 2 has a CPU 3, an I/O interface 4, a ROM (Read Only Memory) 5, and a RAM (Random Access Memory) 6. In the analysis processing section 2, a circuit simulation program is loaded into the RAM 6 via the CPU 3, and the CPU 3 executes it. An execution result is stored in the storage unit 1 via the I/O interface 4, and is outputted to the calculation result outputting unit 7. The calculation result outputting unit 7 is exemplified by a display, a printer, and a file system, and is connected to the I/O interface 4, and outputs the analysis result of the analysis processing section 2.
Next, an operation method of the conventional circuit simulation system will be described. FIG. 1 is a flowchart showing the conventional circuit simulation method.
First, the circuit simulation is started (Step S100). The circuit data 10 and the analysis condition data 20 are read from the storage unit 1 (Step S101). As described above, the circuit data specifies the configuration of the circuit to be analyzed, and contains data indicating the elements in the circuit and connections of them. The analysis condition data is the setting data for performing analysis, which includes a data indicating the tolerance of accuracy.
Next, an initial value set of the Newton iteration method is set. Moreover, a counter variable i indicates the number of times of the iteration of the Newton iteration method and is set to "1" (Step S102).
A circuit equation is generated based on the circuit data 10, and an approximating process of calculating a solution is executed in which the circuit equation is approximately solved by the Newton iteration method (Step S 103). The process of calculating the solution of the circuit equation is matrix calculation. As the circuit scale becomes larger, the matrix scale becomes larger and the number of times of the operation increases exponentially. Specifically, the number of times of the calculation is the number of times of the updating calculation when erasing variables in a LU (Lower/Upper) decomposition of the circuit matrix into a product of triangular matrixes. It is generally known that the number of times of the calculation increases in proportion to the third power of a matrix size.
It is determined whether the solution of the circuit equation in the approximating process has been converged through the iterative calculation (Step S104). Conditional branching is performed based on a result of convergence determining process (Step S105).
When it is determined that the solution has converged (Yes at Step S105), the solution of the circuit equation is saved as the waveform data 40 in the storage unit 1 (Step S108). Then, the operation follow is completed (Step S109).
When the solution has not converged (No at Step S105), it is determined whether the counter variable i has exceeded a predetermined number of times of the iteration ITL1 (Step S106). When it exceeds an upper limit (Yes at Step S106), it is determined to be in error and the operation flow is completed (Step S107). When it does not exceed the upper limit (No at Step S106), the Step S113 is carried out.
At Step S113, the counter variable i is incremented by "1" (Step S113), and the operation flow returns again to the Step S103. The values calculated at the last Step S103 are used as the initial value set, and the process of calculating a solution is executed, like the last time. Therefore, the number of times of the operation required to calculate the solution can be expressed as follows:(Number of times of the operation required to calculate a solution)=(Number of times of the operation per iteration of Newton iteration method)×(Number of times of the iteration required until convergence)The number of times of the operation per iteration of the Newton iteration method can be found at the Step S103, and the number of times of the iteration of the Newton iteration method required until convergence can be found as a value of the counter variable i at the Step S108. Thus, the number of times of the operation required to calculate the solution can be expressed as a product of these numbers of times. A simulation time is dependent on the number of times of the operation required to calculate the solution.
FIG. 2 is a flowchart showing a convergence determining process of the Step S104 of FIG. 1 in detail. First, when the Step S104 is started (Step S500), a convergence variable flag eflg is initialized to zero (Step S509). The convergence variable flag eflg indicates whether there is a variable that has not converged, among variables of the unknown variable X. When the convergence variable flag eflg is "1," it indicates that one or more variables included in the unknown variable X are not in a state of non-convergence; and when the convergence variable flag eflg is "0," it indicates that all the variables are in a state of convergence.
Next, the following process is performed on all the variables included in the unknown variable X. It is evaluated over all the variables whether a difference between the present solution x(i) of the variable x included in the unknown variable X and a last time solution x(i−1) thereof, is within a tolerance (Step S501). When the solution has not converged, it is determined that there is a variable not having converged in the variable X, and the convergence variable flag eflg is set to "1" (Step S507). After that, it is determined whether there is an undetermined variable (Step S502). When there is the undetermined variable (Yes at Step S502), the next variable is evaluated (Step S508). By repeating this, the convergence determining process is made in sequence on all the variables included in the unknown variable X.
When there is no undetermined variable (No at Step S502), the convergence determining process is completed (Step S504).
FIG. 3 is a table showing a convergence determination result. In FIG. 3, "i" represents a value of the counter variable of the Newton iteration method. Regarding "V" and "I," when a set of voltage variables included in a set of unknown variables X is specified as V (V=(xv|xv is a voltage variable, and X includes xv)) and a set of current variables included in the set of unknown variables X is specified as I (I=(xi|xi is a voltage variable, and X includes xi)). A symbol "O" indicates a case that all the elements in each set have converged, and a symbol "X" indicates a case that any of the elements therein has not converged. Here, X is a logical OR of V and I. A column of "determination" shows the convergence determination result of Step S104.
First, at the Step S103 of FIG. 1, when the counter variable i of the Newton iteration method varies to a value not exceeding an upper limit ITL1 of the number of times of the iteration, i.e., a value (n−1) (n: natural number), the variable sets V and I result in non-convergence. At this time, when the counter variable i is incremented to n, namely when the nth iteration process is executed, the sets V and I of FIG. 3 result in non-convergence. At the Step S103, the generation of the circuit equation and the process of calculating the solution by the Newton iteration method by using the (n−1)th solution as an initial value set are performed.
Next, at the Step S104, it is determined whether a value of Xn−X(n−1) stays in the tolerance, that is, whether a difference between the present solution x(i) and the last solution x(i−1) is within the tolerance for the variable x included in the unknown variable X. Here, in FIG. 3, the set V has converged, but the set I has not converged. Therefore, it is determined that the result of the Step S104 is non-convergence. For this reason, the iteration process for the (n+1)th iteration and succeeding iterations are carried out again. In this way, the Newton iteration method is carried out repeatedly until the solutions included in the set V and the set I are all converged, and the number of times of the iteration required until calculating the solution becomes m (n<m).
The process of calculating the solution at the Step S103 is repeated due to the result of the convergence determining process at the Step S104. Even when the voltage variable has converged, the process of calculating the solution and the convergence determining process are repeated until all the variables have converged.
As described above, in the conventional circuit simulation method, since a predetermined determining process is executed to a change in the variable and the solution, the number of times of the iteration increases and a circuit simulation time elongates depending on a circuit size.
Therefore, a circuit simulation method is demanded in which the number of times of the operation of the Newton iteration method until convergence is attained can be reduced by changing a method of convergence determination based on the solution of a variable included in an unknown variable X, and also, the number of times of the operation required to calculate the solution can be reduced as a whole, and a circuit simulation system therefore.

System and method for circuit simulation


An electrostatic spray module for applying agricultural liquids such as a pesticide to crops where, externally to the spray module the number of connections is reduced to three, one for the liquid pesticide, one for compressed air and one for a low voltage signal. Internally to the spray module, the low voltage is converted to a high voltage signal, which is, along with the pesticide and the compressed air delivered to one or more electrostatic spray nozzles using only two electrically conductive pipes, a gas delivery pipe and a liquid delivery pipe. The nozzles fit into the gas delivery pipe and draw the compressed air through gas channel openings in the side of the nozzles. The gas delivery pipe doubles as the means to delivery the high voltage signal to the nozzles. Each nozzle has a liquid feed from the liquid delivery pipe, which carries ground voltage, maintaining the liquid at ground voltage. The grounded liquid merges with the compressed air in the nozzles to form an atomized liquid. The atomized liquid then passes through an electrode, which is electrically charged by the high voltage signal to form an electrostatic spray. The electrical charge in the spray leads to better dispersal of the spray due to the droplets in the spray repelling from each other, and further improves the adherence of the spray to crops which attract the charged droplets.

This invention relates generally to electrostatic sprayers and more particularly to an electrostatic spray module for agricultural applications.
Electrostatic sprayers are commonly used in agricultural applications to apply pesticides and other agricultural chemicals to crops. Electrostatic sprayers typically work based on the following basic principles. Compressed air and liquid (for agricultural applications this is typically a pesticide) are separately piped into a nozzle, where the two mix in the process of atomization, forming small droplets of the liquid. The atomized liquid then passes through a charged electrode, in the process of induction, which charges the droplets. The droplets then, due to the flow of air, spray out to the crops. The charge on the droplets causes the individual droplets to repel from each other scattering the spray for an even and wide spread application. The charge on the droplets also leads to better application by causing the droplets to better adhere to the crops, which are at ground potential and electrically attract the charged droplets.
To sufficiently charge the liquid in an electrostatic sprayer, the nozzle needs a high voltage power source of generally one kilo-volt or higher. Electrostatic sprayers designed for use in the field use a low voltage power supply, such as a 12 V battery (typically the tractor battery), that is hooked up to a power supply that generates a high voltage signal. Some sprayers use one power supply for many nozzles, but have the complications of distributing high voltage over several nozzles. Many sprayers use one power supply per nozzle to keep the high voltage local to the nozzle but have the disadvantage of the high cost and maintenance of the many power supplies.
Systems incorporating individual electrostatic nozzles modules tend to be systems that are cumbersome to hook-up, difficult to configure, hard to maintain, have limitations in their performance and are costly to manufacture.
A "sprayer" in agricultural parlance refers to a electrostatic spray nozzle system or module as well as the supporting components of an air compressor, a liquid tank and pump, a frame and a boom, all of which are typically mounted on a tractor.
A standard sprayer typically comprises 30 to 80 nozzles each with several hoses and connections. Each nozzle in the sprayer requires a hose leading from the compressed air source, a hose leading from the liquid source and either a low voltage or high voltage power supply connection. Multiplying these three connections over the dozens of nozzles that are typically hooked to a tractor boom and the result is a cumbersome process for hooking up a sprayer system. This process can take considerable labor resources and result in a system that is extremely difficult to maintain.
Traditional sprayer systems also tend to have a fixed configuration, making changes to such things as the size of the nozzle opening and the spacing of the nozzles difficult, if not impossible, to alter after manufacturing. Changing the size of the opening in a nozzle affects the airflow through the nozzle, resulting in either lower airflow through the nozzle or a higher pressure. Varying the size of the nozzle opening may be desirable based on the type or stage of a crop. For example, a vineyard in springtime will consist of small plants where it will be preferable to use less air volume or pressure than later in the season when the plants are larger. Changing the spacing of the nozzles can also be difficult in prior art sprayer systems. A fixed distance between nozzles may not be desirable as different crops and different conditions have different nozzle requirements. For example, a sprayer used at a golf course will want full coverage over a flat surface, which will optimally be a nozzle every four inches or so. Meanwhile, a cotton crop has rows spaced such that a nozzle every twelve inches would provide adequate coverage.
Prior art sprayer systems are also difficult to maintain. The prior art nozzle described in Cooper et al., U.S. Pat. No. 5,704,554, shown in FIG. 1, has components, such as the electrode 2 and the liquid channels 5 integrated into the nozzle. Parts such as these may need routine cleaning for optimal use of the nozzle. Build up of dirt and liquid on the electrode can lead to inefficient spraying and excessive current draw on the power supply. Cleaning of the embedded electrode is awkward and can lead to damaging the surface of the electrode and the plastic enclosing the electrode. Cleaning the non-removable liquid channel is difficult. The tip of the liquid channel is subject to build-up of conductive deposits that create electrical current pathways and can lead to carbon deposits building up on the liquid channel tip. Excessive build up can damage the tip of the liquid channel, resulting in an inoperable nozzle. Excessive damage to the electrode or the liquid channel in a nozzle where such parts are non-replaceable requires complete replacement of the full nozzle.
Prior art sprayer systems also have performance limitations. The thin electrode 2 shown in the prior art nozzle of Cooper, et al., in FIG. 1 only provides a limited charge to the droplets, not fully maximizing their ability to attract to crops. Many prior art nozzles also have problems with the flow of the atomized liquid through the tip of the nozzle. This is caused by imperfections in the inner orifice wall of the tip of the nozzle. For example, the prior art nozzle of FIG. 1 has a stainless-steel electrode 2 embedded between plastic 1 and ceramic 3 layers, resulting in a three-layered passage from the end of the liquid channel 5 to the opening after the electrode 6. This three-layered channel of dissimilar materials, even with quality machining, has microscopic notches between the layers of materials. These notches magnify with wear and tear on the nozzle and the wear and tear, in turn, is accelerated by the damaging effects caused by the notching. As the air and liquid mixture flows out of the channel, the mixture eddies along the notches resulting in decreased charging of the spray, increased current draw and physical deterioration of the inner surface of the nozzle. The notching and the dissimilar materials can cause the liquid to be deflected to the side in its passage through the nozzle, causing an inefficient spray pattern and sub-optimal charging by the electrode. Another performance limitation in some prior art nozzles is an off-center spray, resulting from liquid channels that are not in complete coaxial alignment with the output of the nozzle. FIG. 1 shows a prior art nozzle that has a liquid channel with one end of the liquid channel centered about the electrode, but not the entire channel at the end is not straight upstream from the opening. An off-center liquid channel can result in a spray that is not centered around the end of the nozzle due to a lateral force in the liquid generated in the liquid's passage through the off-center liquid channel 5. Further, this can lead to plugging of the nozzle. Prior art nozzles also passively rely on the liquid maintaining ground potential which leads to unreliable charging of the spray.
Lastly, prior art sprayer systems are expensive. Due to the high cost of sprayer systems, their use is generally limited to high-value cash crops and specialty applications such as vineyards. The expense of prior art sprayer systems limits their use in commodity crops.
Therefore, there is a need for a system for electrostatically spraying agricultural crops that is easy for a user to set up, convenient to configure to different field situations and simple to maintain. Further, there is a room for improvement in the performance of sprayer systems as well as a strong need for more affordable systems.


In a semiconductor device, a self-test circuit includes a write part for writing data in a given address of a special region of a nonvolatile memory; a read part for reading the written data from the given address; a verify part for determining whether or not the written data accords with the read data; and a decision part for determining soundness of the nonvolatile memory on the basis of a result of determination made by the verify part. In the case where the written data accords with the read data, the decision part determines that the nonvolatile memory is sound, and in the case where the data do not accord with each other, it determines that the nonvolatile memory is unsound.

The present invention relates to self-test technique for a nonvolatile memory, and more particularly, it relates to memory self-test method and circuit for checking the soundness of a nonvolatile memory necessary to be in sufficient security.
In an LSI for storing highly confidential data or programs, or an LSI including a security circuit for an IC card or the like, significant data or programs with high confidentiality are stored in a nonvolatile memory such as a flash memory or a FeRAM. For example, in a nonvolatile memory of an LSI for an IC card, personal information, authentication data or cryptographic key data for a cryptographic system is stored. Since very significant contents are thus stored in a nonvolatile memory, it is necessary to provide means for preventing a malicious third party from decoding or modifying the data. For this purpose, it is effective to encode data to be stored in a nonvolatile memory.
In order to further increase the security level, information corresponding to a security status is stored in a nonvolatile memory so as to restrict an invalid access itself. In some conventional technique for security protection, the number of fails made in authentication for accessing a nonvolatile memory is written in the nonvolatile memory, and when the number exceeds a given value, it is regarded that an invalid access is being made. Thus, authentication to be made thereafter is rejected, confidential data stored in the nonvolatile memory is deleted, or the starting time in next power supply is elongated.
Even when such means for protecting the security is provided, however, an invalid access cannot be effectively restricted unless the number of authentication fails is correctly incremented. In the case where physical modification for disabling data write is made on a nonvolatile memory by, for example, fixing a write enable signal to a disable state or disconnecting a power line necessary for data write, information corresponding to a security status stored in the nonvolatile memory cannot be updated. As a result, even when authentication fails are repeatedly made, the number of fails is not incremented, and hence, it is impossible to restrict the invalid access.
It can be checked whether or not a nonvolatile memory has been physically modified, namely, the soundness of the nonvolatile memory can be checked, by verifying written data. However, when the verification is performed every time data is written, the power consumption is increased and the data write speed is lowered, and in addition, the lifetime of the nonvolatile memory may be shortened because the number of memory accesses is thus increased.

Memory self-test circuit, semiconductor device and IC card including the same, and memory self-test method


The invention makes it possible to measure binding of a biochemical substance with a high throughput and with high sensitivity using a small cell capable of being filled with a small amount of chemical solution. A space between a first substrate and a second substrate such that probes are immobilized on their mutually facing planes is used as a cell that houses a specimen solution. Light is irradiated from a first substrate side, and reflected light is subjected to spectroscopy. Binding of the target with the probe is detected by a wavelength shift in the reflection spectrum.

Conventionally, measurement of binding between biochemical substances, such as an antigen-antibody reaction, has been generally performed by using labels, such as radioactive substances and fluorescent substances. The use of labels needs time. Especially, in using them in a protein, there are a case where the method is complicated and a case where properties of a protein will change.
To circumvent these problems, a biochemical sensor that uses a change of interference color of an optical thin film is known as a method of directly measuring the binding between biochemical substances without using a label. A paper by T. Sandstrom, et al., Applied Optics, 1985, 24, 472-479 (nonpatent document 1) describes this biochemical sensor. Its example will be explained using a model of FIG. 1. An optical thin film 1-2 is formed on a substrate 1-1. The refractive index of air is 1.00, and the optical thin film 1-2 is a material of a refractive index of 1.50. The substrate having a refractive index of 2.25 is used. If an optical thickness of the optical thin film is chosen to be ¼ of a visible light wavelength λ0 or one of its odd multiples (¾λ0, 5/4λ0, etc.), the optical thin film acts as an antireflective film, producing an interference color. On this optical thin film 1-2, a monomolecular layer of a first biochemical substance 1-3 is formed. If the biochemical substance is considered a protein, its refractive index is of the order of 1.5 and its layer thickness is of the order of 10 nm. At this time, as shown by a reflection spectrum A of FIG. 2, the intensity of reflected light in a direction perpendicular to the optical thin film becomes zero at wavelength λ0. When a second biochemical substance 1-4 forms bond with this first biochemical substance 1-3 biochemically, a change in the reflection spectrum from a solid line A of FIG. 2 to a dashed line A′ occurs, causing the interference color to change. By this change, binding of the second biochemical substance is detected. As a general procedure of detection, first, the optical thin film 1-2 on the substrate 1-1 covered with the monomolecular layer 1-3 of the first biochemical substance is prepared. This is immersed in a solution of the second biochemical substance. Subsequently, it is taken out from the solution and dried, a change of the interference color from the solid line A of FIG. 2 to the dashed line A′ is examined. Moreover, the document describes that the use of a material that is an optical absorbing material, for example, silicon, as a material of the substrate 1-1 can suppress an effect on the measurement caused by the optical reflection generated by the back of the substrate. As in the above, the nonpatent document 1 describes a technique wherein, after the sensor is taken out into air and dried, the interference color is measured.
On the other hand, if a material of a refractive index of approximately 2.2 is used as the optical thin film, a clear interference color can be obtained in an aqueous solution, and accordingly the amount of binding of the first biochemical substance and the second biochemical substance can be measured in real time in the aqueous solution (see a paper by T. Fujimura, et al., Jpn. J. Appl. Phys. 2005, 44, 2849-2853; nonpatent document 2). Its example will be explained using a model of FIG. 3. An optical thin film 3-2 is formed on a silicon substrate 3-1. The optical thin film 3-2 is a material of a refractive index of 2.2 and its thickness is specified to be 70 nm. On this optical thin film 3-2, a monomolecular layer 3-3 of the first biochemical substance is formed. White light is made incident on that structure through an optical window 3-4 made of a transparent material, and a reflection spectrum of the sensor is measured. Moreover, if a bundle of optical fiber is used as a light guide for irradiating white light and collecting reflected light, the size of the sensor can be designed to be of a diameter of submillimeter. FIG. 4 shows the reflection spectrum. In calculation of the reflection spectrum shown in this FIG. 4, since reflection of the light on the surface of the optical window 3-4 hardly affects the measurement, it is ignored. This is done because, by setting a separation between the optical window 3-4 and the optical thin film 3-2 to, for example, approximately 0.15 mm, optical interference between the optical window 3-4 and the optical thin film 3-2 can be prevented from affecting the measurement. The refractive index of a material between the optical window 3-4 and the optical thin film 3-2 was set to a refractive index of water, i.e., 1.333. A layer 3-8 of the first biochemical substance and a layer 3-5 of the second biochemical substance are both specified to be a layer of a refractive index 1.5 and a thickness of 10 nm. A solid line B of FIG. 4 shows a reflection spectrum in the case of absence of the second biochemical substance layer 3-5; a dashed line B′ of FIG. 4 shows a reflection spectrum in the case of presence of the second biochemical substance layer 3-5. If the second biochemical substance forms bond with the first biochemical substance, a change from the solid line B to the dashed line B′ will occur and a minimum position of the reflectance will move to a longer wavelength side by 13.5 nm. By measuring this change, the binding of the second biochemical substance with the first biochemical substance can be measured. Here, in the nonpatent document 1, the refractive index of the layer of a biochemical substance is set to 1.5, and it can be estimated that the layer of a biochemical substance material 3 nm thick within a dimensional range of 10 μm×10 μm contains an organic material of 0.5 pg. From the estimate of this nonpatent document 1, a change of the minimum position of the reflectance of 1 nm in the nonpatent document 2 can be approximated to the amount of the binding of the biochemical substance of about 1 ng/mm2. Since, by this method, measurement of binding can be done in real time in a specimen solution, saturation of a reaction can be found without taking out the sensor from the specimen solution; therefore, it can perform measurement more correctly and more quickly than the method of the nonpatent document 1.

Cells for biochemical analysis, kit for biochemical analysis, and biochemical analyzer


Provided is a therapeutic agent capable of treating allergic rhinitis in a safe, simple, and noninvasive manner substantially without side effects and physical distress to patients. Provided is a therapeutic agent for allergic rhinitis for 5-aminolevulinic acid-based photodynamic therapy (ALA-PDT) comprising 5-aminolevulinic acid (ALA) or a derivative thereof, or a salt of the 5-aminolevulinic acid or the derivative as an active ingredient, and used in ALA-PDT in which light having a wavelength of 400 nm to 700 nm is irradiated; and particularly provided is a therapeutic agent for allergic rhinitis in a locally applied solution form, a water-soluble ointment-dissolved form, a jelly-dissolved form, or the like. The therapeutic agent does not require the conduct of 5-aminolevulinic acid-based photodynamic diagnosis (ALA-PDD) for detecting a site of protoporphyrin IX (PpIX) accumulation radiating light having a wavelength of 610 nm to 650 nm under irradiation of excitation light having a wavelength of 380 nm to 420 nm before ALA-PDT.

Allergic rhinitis is caused by the response of the autoimmune system in the body to a causative agent in the external environment, and dust, mold, pollen, or the like is known as the causative agent.
A drug therapy, a desensitization therapy, an operative therapy by inferior concha cauterization using a laser, or the like is known as a therapeutic method for allergic rhinitis represented by pollinosis. For the drug therapy, a steroidal anti-inflammatory drug, an immunosuppressive drug, an antihistamine drug, or the like is used; however, this therapy is known to have problems in terms of side effects and efficacy. For example, it is considered problematic in that the steroidal anti-inflammatory drug has side effects such as adrenal atrophy and dysfunction and gastric ulcer (see, for example, Non-patent Document 1), that the immunosuppressive drug has side effects such as infection (see, for example, Non-patent Document 2), and that the antihistamine drug has side effects such as malaise, drowsiness, and dizziness. The antihistamine drug is also known to be unsatisfactory in terms of efficacy compared to the steroidal anti-inflammatory drug and the immunosuppressive drug.
The desensitization therapy is known as a method for suppressing only a concerned antigen-specific immune response by specifying an antigen responsible for allergy and intradermally administering the antigen to induce desensitization; however, it is considered to take several months to several years to exert a sufficient effect. In addition, this therapy is known to have a risk of inducing anaphylactic shock or the like due to the direct systemic administration of an antigen causing allergy, making it necessary to gradually increase its dose from a low dose. In other words, the desensitization therapy is considered problematic in that it not only requires a long-term treatment but also is accompanied by a patient's distress due to injection.
The operative therapy by inferior concha cauterization using a laser is considered problematic in that it not only inflicts a physical distress on patients with allergic rhinitis but also is not a method for radically treating allergic rhinitis.
On the other hand, recently, PDT has been developed that is a method involving administering a compound responsive to light and treating a target area by light irradiation. PDT is simple in treatment, exhibits low invasiveness into a living body, and enables the conservation of an organ, for example; thus, it has recently been a focus of attention as a new method for treating cancer in which consideration is given to quality of life (QOL).
ALA as one of the drugs used for PDT is known as an intermediate in a pathway of pigment biosynthesis broadly present in animals, plants, and fungi, and is typically biosynthesized from succinyl-CoA and glycine by 5-aminolevulinic acid synthetase. Although ALA has no photosensitivity per se, it is known that ALA is metabolically activated to protoporphyrin IX (hereinafter also referred to as "PpIX") by a series of enzymes of the heme biosynthetic pathway in cells and then specifically directly accumulates in tumor tissue and neovascular vessels, and that, when laser light is irradiated on the site in which PpIX accumulates, cancer cells are degenerate/necrotize by reactive oxygen species generated by the excitation of light, which are represented by singlet oxygen, hydroxyl radical, superoxide, and the like.
Since Professor Kennedy of Queen's University, Canada, reported in 1986 that skin cancer can be treated by applying ALA and irradiating light (see, for example, Non-patent Document 3), methods have been reported for diagnosing and treating lesioned parts of various sites using ALA. For example, tumor diagnosing agents are proposed which have been developed based on the finding that, when ALAs is administered into the body, PpIX derived from the ALAs accumulates in cancer and emits fluorescence by light irradiation (see, for example, Patent Documents 1 and 2).
Because PpIX emits red fluorescence having a peak at a wavelength of 636 nm when receiving excitation light having a wavelength around 410 nm, ALAs are used for the diagnosis of tumor by 5-aminolevulinic acid-based photodynamic diagnosis (ALA-PDD), and is also expected to be used in applications such as the diagnosis of brain tumor or bladder cancer and the prevention of anemia.


A game including multiple playing pieces; with the multiple playing pieces including a plurality of identical sets of playing pieces. Each identical set of playing pieces includes a plurality of playing pieces having non-numerical indicia representing inanimate or animate objects which are rankable in a specific order relative to each other; each of the plurality of playing pieces including indicia representing one of the objects. Each of the plurality of playing pieces includes retaining means for retaining additional indicia that indicates a specific condition of the object represented on that respective playing piece, with the specific condition having a bearing on the order of ranking of the object.

This invention relates generally to games, and more specifically to an educational game which teaches and reinforces information regarding a variety of objects; both animate and inanimate, and which is readily adaptable to be played as a "card" game.
Card games are extremely popular with both adults and children. Many of the games are extremely convenient to play, because they do not involve the use of other objects, such as dials, dice, boards, etc.
One card game which has been popular for years, is the game of "war". In that game a conventional, 52-card deck of playing cards is dealt, with the cards face down, to two or more players. After all of the cards have been dealt, each player turns his or her top card over, exposing the "value" of the card. The player having the highest card value wins the other players' cards.
If two or more players expose cards having the same value, a "war" is declared between (among) those players. To determine the winner of the war, each of the players participating in the war places one or more cards (most commonly two cards) face down, and then an additional card face up. The participating player with the highest additional face-up card wins the war, and if that card also is higher than the face-up cards of the other players, then he or she wins all of the other players' cards in that particular round. If the additional face-up card of a participant in the war either is the same value as the additional face-up card of another participant(s) in the war, or is the same value as one of the other players, then a second war exists. This second war, and any subsequent wars, are "fought" in the same way as described above with respect to the first war. Play proceeds until one player has one all of the playing cards from the other players.
The game of war has achieved such a high degree of popularity that patterning educational games in a format similar to "war" is an effective way of providing an educational experience for the players; both children and adults. It is to such a game that the present invention relates.


An excess flow valve includes: a valve that opens and closes a flow path; a solenoid that keeps the valve in an opened state; a pressure detector detecting the primary-side pressure of the valve; and a controller controlling current flow to the solenoid based on pressure detected by the pressure detector so that the valve is closed when the flow rate supplied to the valve is at or more than a predetermined amount. The excess flow valve can be used in a wide variety of applications and prevents a flow of fluid at or greater than a desired flow rate even if the primary-side supply pressure varies.

Conventionally, when using a pressurized fluid supplied from a high-pressure tank or another source, an excess flow valve is used to prevent the fluid from damaging devices that it operates and to prevent fluid from flowing out if the devices are damaged. This excess flow valve shuts off its flow path when it receives a flow that reaches or exceeds a predetermined flow or when there is an irregular drop in the pressure on the downstream side of the valve.
FIG. 3 shows an example of a conventional excess flow valve. An excess flow valve 100 includes: a valve body 102 that opens and closes a flow path 101; a spring 103 that keeps the valve body 102 in an opened state; and a valve seat 104 on which the valve body 102 is seated. A housing 105 supports the valve body 102 and the spring 103.
The valve body 102 includes: a disc-shaped base end section 106; a cylindrical guide 107 that extends from the base end section 106 to the right in the figure and that has a diameter smaller than that of the base end section 106; an opening 108 formed at the outer perimeter surface of the guide section 107; and multiple projections 109 projected to the left from the base end section 106.
In the normal state of the excess flow valve 100, the biasing force from the spring 103 keeps the valve body 102 abutted against the left-side end surface of a valve body housing section 110 in the housing 105. The fluid passes from the left side in the figure (primary side) to the right side (secondary side) through the gap formed by the projections 109 and the opening 108 of the guide 107.
With this excess flow valve 100, when the fluid from the primary side has a flow that reaches or exceeds a predetermined amount, the force (flow velocity) of the flow causes the valve body 102 to move to the right in the figure against the biasing force from the spring 103 and to become seated on the valve seat 104, thus blocking the flow path 101 and shutting off the excessive flow.
Under normal conditions, the valve body 102 is kept in an open state by the biasing force from the spring 103 with almost no change in pressure between the primary side and secondary side of the valve body 102. If, however, for some reason the pressure on the secondary side drops to or below a predetermined pressure, the valve body 102 moves to the right in the figure in opposition to the biasing force from the spring 103 and is seated on the valve seat 104, thus blocking the flow path 101.
The background technology described above covers general subject matter, and the present applicant has not, at the time of this application, found any particular document that describes this background technology.
However, with the conventional excess flow valve 100, if the biasing force from the spring 103 stays constant, the flow of fluid that would activate and close the valve body 102 would vary depending on the pressure on the primary side. In general, with fluids, the relationship between pressure and flow velocity is such that, when the flow rate is constant, an increase in pressure results in a decrease in the flow velocity. With the valve body 102, operation is based on the flow rate, and this flow rate can be expressed as cross-sectional area x (multiplied by) flow velocity. As a result, when the primary-side pressure increases, the flow velocity decreases even if the flow rate is the same and less force acts on the valve body 102, keeping it open. To make the valve body 102 close while keeping the pressure unchanged requires increasing the flow velocity. This means that the higher the primary-side pressure is, the flow must necessarily increase.
Because of this, the conventional excess flow valve 100 requires the selection and mounting of a spring that has a biasing force that closes the valve at the desired flow rate for the particular primary-side pressure to be used. This is inconvenient and prevents the valves from being used in general-purpose devices.
In addition, if the valve is used in something where the primary-side pressure changes, the flow rate that closes the valve body 102 would change depending on the primary-side pressure. This makes it impossible to keep the flow rate consistently below a predetermined value.
If the biasing force of the spring 103 in the excess flow valve 100 is set so that the valve closes at a desired flow rate under a high primary pressure, a drop in the primary-side pressure would result in the valve closing at a flow rate that is lower than the desired flow rate, thus preventing the flow rate needed on the secondary side from being provided.


A strip member is wound by a tire building machine having a building drum for taking up a strip member of predetermined length without discrete ends and a transfer conveyor for feeding the strip member to the building drum. The winding-start end of the strip member transferred from the transfer conveyor to the building drum and the winding-finish end of the strip member are detected by means of end sensors, measuring the angle through which the building drum has rotated from the output of a winding start signal to the output of a winding finish signal in response to output signals from the sensors by means of an angle detector. A lapping margin between the winding start end and the winding finish end of the strip member is calculated based on the rotational angle and diameter of the building drum by means of an operational means.

1. Technical Field
The inventions disclosed in this application relate to methods for winding an automotive tire building strip member on a building drum.
More particularly, the first invention in this application provides a method for winding a strip member, such as rubberized steel cord, tread rubber or the like, precut to the circumferential length of the building drum of a tire building machine, on the building drum in such a manner that both ends of the strip member will be exactly abutted against each other on the drum.
The second invention in this application relates to an automotive tire building device wherein the lapping margin of ends of a strip member wound on a building drum can be automatically measured.
2. Prior Art
In winding a tire strip member on a building drum, it is at times required to insure exact abutment of its ends and at other times required that the lapping margin (positive and negative) be within certain allowable limits.
As a technology for insuring exact abutment of both ends of a strip member, such as rubberized steel cord or tread rubber, in the winding thereof on the building drum of an automotive or other tire building machine, there is known a method described in Japanese Patent Publication No. 61-32980. According to this method, the length of the strip precut to the circumferential length of the building drum is measured and, then, a leading portion, of predetermined length, of the strip is taken up on the building drum at the feeding (takeup) speed equal to the peripheral speed of the drum. Then, the feeding speed relative to the building drum speed is altered to either compress or stretch the strip while its intermediate portion is wound on the building drum. Finally, the feeding speed is re-equalized with the peripheral speed of the building drum to wind up the remaining portion of the strip member on the drum.
However, the strip member tends to shrink with the progress of time after cutting and the time to termination of shrinkage and the amount of shrinkage is dependent on the environment and other conditions. Under certain conditions, the amount of post-cutting shrinkage reaches as much as about 0.5%. Mor over, the length of the strip member varies with the magnitude of the tension that acts on the strip when it is transferred from a transfer conveyer to the building drum and the pressure of contact between the strip and the drum. In the prior art method mentioned above the total length of the strip member is measured while it is undergoing shrinkage after cutting and the ratio of the feeding speed to the peripheral speed of the building drum is set according to the length value thus found so as to compress or stretch the intermediate portion of the strip. As a consequence, the aforesaid ratio is often irrelevant and there occurs an excess or a shortage of compression or stretching of the strip member, thus causing a lapping of the leading and trailing ends of the strip member or a gap therebetween.
Therefore, in the first invention in this application provides a winding method which insures exact abutment of the ends or abutment without an appreciable overlap or gap.
Furthermore, in the process of manufacture of an automotive tire, not only the steel cord and tread rubber mentioned above but a variety of other rubber-based strip members are cut to length and both ends thereof are joined to build single-layer or multi-layer ring or cylindrical members. For example, on the primary building machine, an inner liner (a strip member of rubber) and a plurality of ply cords (rubberized fiber cords or steel cords), all precut to the circumferential length of the building drum, are laid up on the peripheral surface of the drum. In the secondary building machine, two steel belts (rubberized steel cords), a cap ply (rubberized nylon cord), etc. are laid up some of these different strip members are not joined by abutment at ends as described above but are joined by lapping the trailing end over the leading end on the drum. As mentioned above, the strip member has the property to shrink on standing after cutting to length as mentioned above and the amount of this shrinkages varies with changes in environment. Also as mentioned above, the length of the strip member is altered by external forces that act on the strip when it is transferred from the conveyer to the drum. Therefore, even if the strip member has been cut precisely to length, a variation is inevitable in the degree of lapping of both ends. Therefore, it is common practice for the operator to trim off the rubber or cord at the trailing end when the degree of lapping is too large or unwind the strip partway and rewind it with stretching when there is a gap between ends.
However, the conventional tire building machine is not equipped with very effective means for inspecting the end joint of the strip member and the current trend toward automation of tire building cannot completely avoid a risk of products with the surplus or deficiency in end lapping being shipped uncorrected and marketed.
Accordingly the second invention in this application provides a tire building device which is capable of detecting a surplus or deficiency in end lapping while a strip member is wound on a building drum of the primary or second building machine.

Method for winding strips on the building machine


A base for display units that can be made as a product of the paper-transformation industry including a box-like body with a lateral wall that extends between a bottom part and a top part. The base also includes a wall on which rests a display unit selectively mobile between:                a lowered position in which the mobile wall is situated in correspondence with said bottom part, and the base is capable of acting as a container for a display unit, and        a raised position in which the mobile wall is situated in correspondence with the top part of the base and is capable of supporting a display unit in the exhibiting position.

1. Field of the Invention
The present invention relates to bases for display units.
2. Description of the Related Art
The invention was developed with particular attention to its possible use in conjunction with display units used to exhibit foodstuffs such as confectionery at points of sale.
A current practice entails the production of bases of this type in the form of products of the paper-transformation industry. Substantially speaking these are boxes of cardboard or similar material, generally having the shape of a parallelopiped, destined to stand on the ground and provided with an upper plane on which a display unit is placed (generally of the type with a number of shelves) on which the exhibited products are placed.
The purpose of the base is to maintain the display unit and the products that are on it in an easily visible and accessible position: it is considered undesirable for products to be located close to the floor, a position that also makes the operation of removing the products from the display unit inconvenient.
It is equally general practice to send the unit formed of the display unit (usually with the products already situated inside it) and the relative base to the places of exhibition and sale, as a unit already assembled and mounted. For this purpose, the display units mounted on their respective bases are loaded onto pallets so that they can easily be handled in groups of four or more units.
An intrinsic drawback connected with this solution lies in the fact that the base of the display unit is in the form of a body having rather considerable dimensions, but whose internal volume is completely unutilised. To use an expression widely used in the sector—this solution in practice translates into the operation of "transporting air": a considerable part of the volume of the load of each pallet is occupied by cavities that are not utilised.
To avoid or at least to mitigate this drawback, a known and practised solution is to employ display units that can be taken apart, that are sent to the place of exhibition in a folded condition, after which they are mounted and receive the products in a subsequent loading operation of the mounted display unit.
Apart from all other considerations, this solution is not appreciated by those who must mount and load the display units at the place of exhibition. The operation of mounting the display unit may in some cases be fairly complicated for an operator who is not particularly expert. In any case, it is inevitable, as the final phase of setting up, that the products—which are transported to the place of exhibition separately, and thus inefficiently—must be loaded into the display unit.
When the display unit has completed its function, the problem arises of taking the display unit apart and disposing of it, including the base. This operation may be particularly hard above all if the display unit (and the base supporting it) are made of mixed materials, destined to follow different channels in the context of differentiated waste collection.

Base for display units, for instance display units for food products


A semiconductor single crystal lift device which comprises a crucible for melting materials for forming semiconductor single crystals and a radiation screen disposed at an upper portion of the crucible and formed of a reversed-cone-shaped adiabatic tube surrounding a lift zone, the apparatus being adopted for lifting up the semiconductor single crystal from a melt in the crucible, in which the radiation screen is divided into more than three adiabatic members, at least part of the adiabatic members being configured in a detachable fashion so that an adiabatic nature of the radiation screen can be partly altered.

1. Field of the Invention
The present invention relates to a semiconductor single-crystal lift device for producing semiconductor single crystals by using the Czochralski Method (the CZ method).
2. Description of the Related Art
As disclosed in Japanese Patent Application Laid-open No. 7-172971, a radiation shield screen is used in the process of lifting crystals for the purpose of accelerating the lift speed, avoiding pollution of impurities and improving the yield of dislocation-free crystals. The radiation shield screen is constructed from an upper screen and a lower screen. The upper screen consists of a three-layer structure, which has two outer graphite members and an inner adiabatic material made of graphite or ceramic fibers disposed therebetween. The lower screen comprises a one-layer structure and is made of graphite, quartz or fine ceramics.
In a semiconductor single-crystal lift device provided with above-described radiation shield screen, the radiation shield screen is used to shield single crystals from the radiation coming from the melt and crucible made of quartz, and to expedite the cooling of single crystals. In other words, the radiation shield screen is used for heat obstruction. By way of omitting the adiabatic member disposed near the free surface of the melt, the high temperature zone (above the free surface of the melt) of over 1200.degree. C. can be expanded so as to prolong the time required for passing grown crystals through the high temperature zone. To obtain high dielectric strength of oxidation layers formed on semiconductor wafers, it is necessary to prolong to the maximum the time single crystals pass through the high temperature region (at 1000.degree. C. to 1200.degree. C. ). The above requirement is nowadays well-known to those skilled in this art.
However, in the above mentioned prior art process, at temperatures ranging from 1000.degree. C. to 1200.degree. C., single crystals cool quickly, which causes a decline in dielectric strength of oxidation layers formed on single crystals. On the other hand, if the lifting speed is reduced, dielectric strength of oxidation layers formed on single crystals is enhanced, yet the yield of single crystals decreases. In conclusion, to enhance dielectric strength of oxidation layers formed on single crystals, it is essential to reduce the cooling speed of single crystals at temperatures ranging from 1000.degree. C. to 1200.degree. C. Furthermore, to reduce the amount of OSF(Oxygen Stacking Fault) and oxygen precipitation, it is required to accelerate the cooling speed of single crystals at temperatures below 1000.degree. C.

Semiconductor single-crystal lift device


Disclosed are method and apparatus for registering multiple sensors collecting data from multiple objects. Sensor registration is decomposed into a two-step procedure. The first step corrects systematic errors. The second step assigns objects measured by one sensor to objects measured by a second sensor. Systematic errors are corrected by generating the global minimum of a systematic error function. One embodiment for generating the global minimum uses a Continuous Greedy Randomized Adaptive Search Procedure.

The present invention relates generally to methods for sensor registration, and more particularly to sensor registration via global optimization methods.
For measuring parameters in a system, multiple sensors may be deployed. For example, the temperature in a building complex may be monitored by installing thermometers throughout the complex. If the thermometers are read at different times, the temperature within the complex may be mapped as a function of location and time. In this example, three types of sensors are involved: a position sensor, a thermal sensor, and a time sensor. Depending on the system requirements, sensors may have a wide range of complexity. For example, a position sensor may be as simple as a tape measure, or as complex as a laser interferometer. Similarly, a thermal sensor may be as simple as a household thermometer, or as complex as an infrared imaging system.
For a measurement system comprising multiple sensors, a key process is sensor registration, which may be broadly construed as applying correction factors to sets of data measured by more than one sensor. Details of sensor registration are discussed below. One simple example of sensor registration is the synchronization of two clocks. At the same instant, if one clock reads 2:00 pm, and a second clock reads 2:05 pm, there is a 5 minute offset between the two. If the clocks are adjustable, they may be registered by adjusting them until they read the same; that is, the offset is set to zero. Alternatively, the offset may be corrected when the data is processed. If the first clock is used as a reference, 5 minutes are subtracted from the time read by the second clock. A second example of sensor registration is the calibration of two thermometers. At the same temperature, if one thermometer reads 110 degrees C. and a second thermometer reads 113 degrees C., there is a 3 degree offset between the two. If the thermometers have a calibration adjustment, the offset may be set to zero. Alternatively, in processing of the temperature data, if the first thermometer is used as a reference, 3 degrees may be subtracted from the reading of the second thermometer.
Some sensors measure a localized value. For example, a thermometer reads the temperature at the location of the thermometer. Other sensors sample values over a region. For example, an infrared camera measures temperatures over a region within its field of view. With a measurement system comprising a distributed array of infrared cameras, the temperature over an extended region may be mapped. Furthermore, more detailed characterization of the temperature variation over a localized region may be acquired if the field of view of multiple infrared cameras overlap. The temperature may be averaged over the readings from multiple cameras, for example. The process of aggregating data from multiple sensors is referred to as data fusion. Proper data fusion requires sensor registration. In the example of infrared cameras, temperature offsets among the different cameras need to be determined. Also, the relationship between the coordinates of a point measured relative to a specific infrared camera and the coordinates of the same point measured relative to a common reference coordinate system need to be determined. Furthermore, if the temperatures are measured as a function of time, and if times are read from a time stamp relative to an internal clock in each infrared camera, synchronization of the clocks to a common reference is required.
The number of sensors which a measurement system can accommodate is limited in part by the mode of data collection. For example, if temperature is measured by thermometers installed throughout a complex, the number of sensors may be limited by the number of stations that a technician can visit within a fixed period of time. With sensors which can transmit data remotely, however, complex measurement systems comprising large arrays of sensors connected via a network may be constructed. The size of the sensor array depends in part on the number of sensors measuring the same parameter at different locations (e.g., 100 infrared imaging systems placed at different locations) and in part on the number of sensors measuring different parameters at the same location (e.g., sensors to measure position, time, temperature, air pressure, and humidity). Therefore, a coherent method for registering an array of sensors, in which the array comprises an arbitrary number of sensors of arbitrary types, is advantageous. For dynamic measurement systems, sensor registration in near real time is further advantageous.
In many applications, measurements from multiple objects need to be collected. In one of the examples above, the temperature distribution of a building complex is monitored by placing thermometers throughout the complex. Here the temperatures at multiple locations are collected from multiple thermometers. In this instance, the thermometers are installed in fixed locations. For example, thermometer 1 is installed in room 101 on the first floor, and thermometer 2 is installed in room 201 on the second floor. The identities of the thermometers are well-known, and the probability of mis-identifying a thermometer is small. The probability is not zero, however, since a technician may make an error. For example, "thermometer 2, room 301" is erroneously entered into a database.
In other situations, identification of the objects may not be straightforward. This is especially true if the objects are mobile. Consider the tracking of aircraft by radar systems. For a commercial airliner, the pilot is in contact with air traffic controllers. The airliner broadcasts a unique identification number, and the identity of the plane is maintained as it is passed from one control tower to another. In a military situation, however, the identities of all the planes will probably not be known, especially for enemy aircraft. Distinguishing friendly aircraft from hostile aircraft is, of course, crucial.
In a complex measurement system, then, multiple sets of data are collected from multiple objects by multiple sensors. For correct characterization of the objects, the identification of the objects being measured by more than one sensor must be determined. For a single object being measured by more than one sensor, sensor registration errors need to be corrected. One approach to sensor registration between two sensors involves minimizing a likelihood function associating the measurements of the objects by two sensors. With this approach, sensor registration falls into the category of global optimization problems, which arise in a variety of science and engineering fields, as well as in other fields such as economics.
Existing methods for solving global optimization problems have various limitations, such as, requiring a good initial estimate of the optima or advance knowledge of the gradient of the function to be optimized. Also, some methods may converge on local minima instead of the global minimum. A coherent method for solving global optimization problems which do not have the above limitations are advantageous. Methods which are computationally efficient are further advantageous.

Sensor registration by global optimization procedures

University of Florida Research Foundation, Inc.