Disclosed herein are two separate processes that do not require a propellant and do not produce an equal and opposite reaction against any external form of matter in the Local Inertial Reference Frame and do not violate Newton's Laws in the Universal Reference Frame. The first process produces horizontal motion, relies on the earth's gravitational field as an external force, and has been successfully tested. The second process produces vertical motion and relies only on the aether. It has been successfully tested considering the effect of the earth's gravity. Due to the law of conservation of angular momentum, the first process is considered by some to not be possible, but with the proper use of an external field (for example, gravity) and the phenomenon of precession, it is clearly possible. A clear distinction is made between a simple rotor and a gyroscope which is a far more complex device.

For almost a century there have been close to a hundred patents issued that claim to produce inertial propulsion, usually in the form of converting rotary motion to unidirectional linear motion. NASA funded a program titled "Breakthrough Propulsion Physics" (BPP) from 1996 to 2002. It was a very successful program in that it provided an opportunity for anyone who believed they had a propulsion breakthrough to present their concept. Its goal was to seek the ultimate solutions to the following three main problems: no propellant required, speeds approaching that of light, and a source of energy to power any such devices (for example, zero-point energy). Terms like "Space drives," "Warp drives," and "Wormholes" are now being used routinely and are written about regularly in reputable scientific journals providing a very healthy atmosphere for creative breakthroughs, thanks to the BPP project.
The BPP Project was a success in that it produced 14 peer-reviewed articles. The project was terminated in 2002 due primarily to a lack of funding, but also due to the realization that out of thousands of submissions, nobody had submitted an idea that appeared to work. Many of the submissions to the BPP Project involved concepts that were already known to not work. Most of the concepts were divided into three common categories: Oscillation Thrusters, Gyroscopic Antigravity, and Electrostatic Antigravity. A detailed analysis was given of at least one example in each of the three categories. The analysis would give a description of the device, then state why it looked like a breakthrough, give a reflexive objection as to why the device cannot work, a deeper assessment, a conclusion, and a "What If" in case someone actually figured out a way to make it work. As an example of an oscillation thruster, the "Dean Drive" described in U.S. Pat. Nos. 2,886,9761 and 3,182,5178 was given. As an example of a Gyroscopic Antigravity device, Dr. Eric Laithwaite's4 work was mentioned. As an example of an electrostatic antigravity device, various Biefeld-Brown effect devices were mentioned, including Lifters, and Asymmetrical Capacitors.
An example of what appears, at first glance, to be propellantless propulsion, but in reality cannot move its center of gravity, is given by U.S. Pat. No. 5,280,864 and described in reference 8.1.
The first instance known to this inventor of a successful demonstration of movement of a device involving the use of gravity as an external force was made by Alexander Charles Jones on May 20, 1975. Alex Jones demonstrated successful inertial propulsion one cycle at a time to Dr. Eric R. Laithwaite. Alex Jones (now deceased) may be considered as the Father of Inertial Propulsion. A reenactment of this first demonstration may be observed by watching the British Broadcasting Company's video titled the "Heretics". Alex Jones' first patent application was in German and was titled, "Vortriebsvorrichtung" (Forward Thrust Device6), Patent # 23 41 245. The patent was filed on Aug. 16, 1973 and was issued on May 22, 1975. The principal inventor of this current Inertial Propulsion System (H. Fiala) translated the original Jones patent from German to English.
By far the most comprehensive patent to date on the subject of inertial propulsion is that by Dr. Eric Robert Laithwaite (also now deceased), U.S. Pat. No. 5,860,317 titled, "Propulsion System4", filed on May 5, 1995 and issued on Jan. 19, 1999.
The problem with existing space vehicle propulsion systems is that they require large amounts of highly explosive propellants, as can be recalled from the explosion of the Space Shuttle Challenger in 1986 and the explosion of very many rockets on the launch pad for both the United States and foreign countries. Rockets using solid or liquid propellants are clearly a brute force and very dangerous approach to manned space flights and space travel. Zero-point energy9,10,12,13,14,15,16,17 is now recognized as existing even though man has not yet managed to successfully harness it. However, it is anticipated that within a few decades, assuming that zero-point energy will have been developed, combining a zero-point energy source with an inertial propulsion system will constitute a perfect marriage of the two technologies for future travel to the'planets and the stars.
Physics of Inertial Propulsion Employing Sustained Acceleration
Einstein's Special Theory of Relativity states that nothing can travel faster than the speed of light. All of the tests of the Theory of Relativity are subject to observations that are usually made by perceiving the results with radar or visually or with photosensitive devices. It is absolutely true that nothing can be observed to travel faster than the speed of light (FTL) because electromagnetic waves are the medium used to make the observations. If an object is traveling away from an observer at greater than the speed of light, its speed could never be measured using light originating at the observer or being emitted by or reflected off the object. If an object were traveling at greater than the speed of light, then according to Special Relativity its mass would be imaginary as shown by the following equations. Trying to measure an object traveling at a speed greater than that of light using some form of electromagnetic waves is like trying to measure the speed of a B1 bomber flying at Mach 3 (three times the speed of sound) using only sonar. It can't be done. According to Special Relativity, the mass m of an object; as a function of its velocity and the speed of light is:m=m0/√(1−v2/c2) If the velocity exceeded the speed of light by an amount x, thenv=c+x, m=m0/√(1−(c+x)2/c2)=m0/√(1−(c2+2cx+x2)/c2)m=m0/√(1−1−2x/c−x2/c2)=m0/√(−2x/c−x2/c2))=−jm0/(2x/c+x2/c2)m=−jm0/√(2x/c+x2/c2)=−jm0/√((x/c)(2+x/c)
This result is the author's own theory. It is just as believable for the mass to become imaginary as it is for mass to turn infinitely large as its speed approaches that of light. This paragraph may generate some controversy. The author welcomes any physical proof to the contrary. If a light beam or a burst of electrons is split in two with each half going in opposite directions, what is the speed of one wavefront with respect to the other wavefront? It is 2c (twice the speed of light). Examine some one-way cases:
If the velocity were to equal the speed of light, then m = m0/√(1 − v2/c2) = m0/√(1 − 1/1) = m0/√0 = ∞If the velocity were twice the speed of light, then m = m0/√(1 − 4) = m0/√(−3)) = m0/−1.732i = −.577jm0At twice the speed of light, the mass is reduced to its original massIf v = √2, then m = m0/√(1 − 2) = m0/(√(−1) = m0/i = −jm01/i = −I1/j = −jAt 1.414 times the speed of light, the mass is reduced to its original mass.
If the velocity were equal to ten times the speed of light, then m=m0/√(1−100)=m0√/(−99)=−jm0/(9.95)
If the velocity were equal to 100 times the speed of light, then m=m0/√(1−10000)=m0/(−9999)≈−jm0/(100)
If the velocity were equal to 1000 times the speed of light, then m=m0/√(1−106)≈−jm0/(1000)
If can be seen that if the velocity increases to n times the speed of light, the mass goes down by a factor of n where n2>>1.
mv<c = m0/√(1 − v2/c2) = m0/√(1 − n2)for v < c(n < 1)mv=c = ∞for v = c(n = 1)mv>c = m0/√(1 − v2/c2) = m0/√(1 − n2) = −jm0/√(n2 − 1) ≈ −jm0/nfor v > c(n > 1)m = mv<c + mv>c = m0/√(1 − v2/c2) − jm0/√(v2/c2 − 1)(vector form of mass)
The preceding equation is a vector form of mass for (0>v>c). This is no different than the vector forms for voltage, current, and impedance. Whereas Z(impedance)=R(resistance)+jX(reactance), where jX is the imaginary component of the impedance, either capacitive reactance (−jX) or inductive reactance (+jX). For electrical engineering, the letter j is used to indicate the reactance (imaginary) component.
Kinetic Energy (KE) versus the velocity of the starship is the product of its mass times its velocity squared.
m = m0/√(1 − v2/c2) = m0/√(1 − v2/c2) = m0/√(1 − n2) = −jm0/√(n2 − 1) ≈ −jm0/nfor n2 >> 1KE = mv2 = m0 v2/√(1 − v2/c2) = m0 n2c2/√(1 − n2) = −jm0 n2c2/√(n2 − 1) ≈ −jm0 n2c2/√(n2 − 1)) ≈ −jm0 n2c2/nKE ≈ −jm0 n2c2/n ≈ −jm0 nc2 ≈ −jm0 (v/c)c2 ≈ −jm0 vcfor n2 >> 1
Summarizing the expressions for kinetic energy below, at, and above the speed of light,
KEv<c = m0 v2/√(1 − v2/c2) = m0 v2/√(1 − n2)for v < c(n < 1)KEv=c = ∞for v = c(n = 1)KEv>c = m0 v2/√(1 − v2/c2) = −jm0 v2/√(v2/c2 − 1) = m0 v2/√(1 − n2) ≈ −jnm0 c2for v > c(n > 1)KE = KEv<c + KE v>c = m0 v2/√(1 − v2/c2) − jm0 v2/√(v2/c2 − 1)vector form of kinetic energy
This is Einstein's own theory in engineering terminology. At warp speeds, it can be seen that if the speed is n times the speed of light, for large n, the imaginary mass is equal to the original mass divided by n; that is the mass goes down significantly, being inversely proportional to its velocity. If the speed were 10,000 times the speed of light, the mass would go down to only one-ten thousand of the original mass. That should be extremely good news for warp drive technology. However, the bad news is that large n the kinetic energy (KE) goes up proportional to the velocity times m0c2. Where have we seen the term mc2 before? It is Einstein's equation for the energy contained in a mass m.
It is very interesting that after reaching the speed of light the kinetic energy drops down very sharply for very small increases in the speed. It then bottoms out to a value of −2jm0c2 at a speed of the square root of '2' (√2) times the speed of light (c). The expression for the kinetic energy for speeds above the speed of light is KEv>c=−jm0v2/√(v2/c2−1). Setting equal to zero the first derivative of the KE with respect to velocity will determine any points on the curve where the slope is zero. That will determine any maximums or minimums. To differentiate an expression of the form of that for the KE, use the formula (du/dx)[u/y]=[ydu/dx−udy/dx]/y2. Let u=v2 and y=√(v2/c2−1). Setting the first derivative of the KE=0, the solution is v=c√2. This establishes the minimum value of the kinetic energy at v=c√2.
Since the KE has a low point at v=√2, the curve at that point is concave upward. Since the curve will have to reverse its curvature to become asymptotic at v=∞, there will have to be a point of inflection at which the curve changes from concave upward to concave downward. The point of inflection can be obtained by taking the second derivative of the velocity and setting it equal to zero. The result is that the curve changes from concave upward to concave downward at v=2.135c. A study of the kinetic energy curve should provide warp drive theorists and designers to attempt to jump directly from just below the speed of light to about 40% above the speed of light to avoid the problems encountered at the speed of light. One theoretical way to accomplish that is to change the point of reference from which the velocity of the starship is determined to a reference frame that is moving at about 45% of the speed of light in the direction of the starship travel. How to accomplish that is left as an exercise for the student.
I think the biggest bather to developing a warp drive is psychological, with the mistaken notion that the mass remains infinite after passing through the speed of light, but that is not true; it goes down inversely as the velocity increases. Maybe it is when a worm hole is entered and the speed increases beyond the speed of light, that mass diminishes and becomes imaginary. In electrical engineering, voltage, current, and impedance, all have an imaginary components, designated with a "j" ahead of their magnitude. In mathematics, the square root of a quantity may be designated by the familiar square root symbol, but is commonly referred as a 'radical'.
The inventor believes the expressions for the kinetic energy of a starship (−jnm0c2) indicates very strongly that to achieve warp speeds, the energy has to be drawn directly from the vacuum or zero point energy as indicated by the m0c2 term. The stress-energy tensor of the aether mathematically defines each point in four-dimensional space-time28. It is linked to the mass and energy distribution of space and would include the zero-point energy at each point in space. To get to a higher velocity, the starship has to either use zero-point energy directly from the aether or collect enough interstellar hydrogen29 with a hydrogen ramjet and convert it to pure energy according to the equation e=mc2. Once a starship has a very high velocity, it can collect more interstellar hydrogen in a given amount of time. If it increases its speed by a factor of 100, it can collect 100 times as much hydrogen each second and that will allow it to increase its speed still further. It is an exponentially easier process to achieve increasingly greater speeds. The main problem is for the engineers and scientists to get serious about harnessing zero-point energy and heating hydrogen to a high enough temperature to fuse hydrogen into helium and release its excess thermonuclear energy. The astronomers need to map out all the nearest stars and estimate which one has the most hydrogen in the path from our solar system to it. Start the voyage with on-board hydrogen to get up to a high enough speed and then when that weight is shed, the mass will decrease and the acceleration will increase.
The ten nearest stars to our solar system are: Proxima Centauri at 4.2421 light years (LY), Alpha Centauri A and B at 4.3651 LY, Barnard's Star at 5.9630 LY, Wolf 359 at 7.7825 LY, Lelande 21185 at 8.2905. Sirius A and B at 8.5828 LY, Luyten A and B at 8.7280 LY, Ross 154 at 9.6813, Ross 248 at 10.322 LY, Epsilon Iridani at 10.522 LY, and Lacaille 9352 at 10.742 LY.
A Trip to Alpha Centauri A
Alpha Centauri A is the second closest star about 4.3 light years away from our Solar System and it is the nearest large star to the Earth. Assume for the sake of a trip to Alpha Centauri A that inertial propulsion providing sustained acceleration is possible, that a space ship can travel faster than the speed of light, that a constant thrust will result in a linear increase of velocity, and that energy for the trip can be obtained from the vacuum or zero point energy. These may be large assumptions, but you will see that they are worth exploring.
The following calculations do not take into account the possible increase of mass with velocity. Assume that a space ship can accelerate continuously at 1.0 g. The speed of light is c=186,000 mps and g is 32.2 ft/sec2.
V=αt=c (speed of light), where V is the velocity, α is the rate of acceleration, t is time, and c is the speed of light. The time required to get to the speed of light is:    t=c/α=30.5 (106) sec=8,472 hrs=353.0 days to get to the speed of light.
Let distance be represented by "S". To travel half-way to Alpha Centauri,S=αt2/2,t2=2S/α=4,130 (1012) sec2     t=64,300,000 sec=744.3 days=2.04 years to get half way to Alpha Centauri. It would take another 2.04 years to decelerate to reach the star for a total one-way time of 4.08 years. At the midpoint of the trip to the star, the velocity would be maximum and would be:V=αt=2.07(109) ft/sec
In terms of the speed of light, at the midpoint,    V=αt=2.03 c, about two times the speed of light.
The total round trip time equals 4.08 years going, 1.00 year in orbit around Alpha Centauri to make observations and collect data, and 4.08 years returning to earth for a total time of 9.16 years.
At an acceleration and deceleration rate of 2.00 g's instead of 1.00 g, the total round trip time would be 2.04+1.00+2.04=5.08 years.
This means that your wife and children would still be home waiting for you. Your children would be 5 years further along in their education.
It is the interpretation and understanding of the inventor that if a space ship had a mechanical (or electric or magnetic or gravitational or nuclear) inertial propulsion unit (IPU) that could sustain acceleration of the whole vehicle at 1 or 2 g's, after 176.5 days at 2 g's the space ship would pass through the speed of light. It is not a matter of its speed being tracked from the Earth, but rather a simple matter of checking the accelerometer to insure that 2 g's of acceleration is maintained and after 176.5 days the space ship will be at the speed of light. Use a fish scale calibrated in pounds and set the acceleration controls so that a one pound mass weighs two pounds and use a wrist watch to tell when 176.5 days are up and you will be at the speed of light, barring any special relativity effects.
As the speed of light is approached, there may not (in the author's opinion) be a physical barrier such as the sound barrier to contend with (unless the aether "wind" has some surprises for us, which it probably will). The mass of the astronauts and their ship will not have become infinite but will reduce as the speed of light is exceeded. At approximately 353 days, the space ship will be traveling at twice the speed of light. At 372.2 days the space ship would be half way and would begin to decelerate at 2 g's.
The point is that if sustained acceleration can be achieved without a propellant, there will be no need to worry about the ability to travel at or greater than the speed of light. Don't create a problem where none may exist. The author cannot prove that this is the case, but he would be disappointed if someone were able to prove him wrong.
Hopefully by the time that sustained acceleration is developed, so too will zero-point energy be developed as a source for power. This would be a perfect marriage of the two technologies.
Examples of Sustained Accelerations Using Known Technology
Rocket Engines with solid or liquid fuel: A brute force and potentially explosive technique. It is a very short term sustained acceleration dropping to zero when the fuel is exhausted. As the fuel is used, the mass decreases and the acceleration increases. It can produce large accelerations based on current technology.
Ion and Plasma engines: RTG powered ion engines have been used many times. The thrust is small, but the thrust can be maintained for years. The ion engines can also be powered with solar energy. Plasma engines can develop considerably more thrust by heating hydrogen gas to a million degrees and ejecting out the end of the rocket and can produce sustained thrust until the supply of hydrogen is exhausted.
Solar Sail: Would require very large sails to produce very small accelerations, but the thrust can be sustained as long as light of sufficient magnitude from one or more stars is available. Beam power from banks of lasers, such as based on the moon, could beam high energy laser beams to push the solar sail.
Ramjet Engines: Atmospheric ramjet engines use oxygen from the atmosphere instead of carrying on-board oxygen. Short duration tests have reached velocities up to mach 15. Ramjets using fusion of hydrogen are still a thing of the future29.
Applications of Gravity
    Dams for power generation    Pendulum to determine time—Grandfather's clock    Foucault pendulum to demonstrate rotation of the Earth on its axis    Gravity creates vortices—flushing a toilet, whirlpools in the water draining through a culvert    Gravity gradiometer    Mass gradient sensor    Holding satellites in orbit by balancing out centrifugal force    Measuring the charge on an electron (Milliken oil drop experiment)    42 minute gravity propelled trip through a straight line hole from any point on Earth to any other point.    Establishing a local vertical    Establishing a geopotential surface as in a lake    Utilization of Lagrangian points    Converting gravitational potential energy by coasting downhill to start a vehicle engine    Horizontal Motion by Mass Transfer (HMT) requiring a gravitational field

Inertial propulsion device to move an object up and down


A computer system and a computer-implemented method for optimizing the number of conditioning data used in multiple point statistics simulation. The method includes inputting a training image representative of subsurface geological heterogeneity; and inputting an initial conservative number of conditioning data. The method further includes selecting a geometry of a template wherein a size of the template is defined by the conservative number of conditioning data; building a search tree using the template by scanning the training image with the template and storing data patterns present in the training image in the search tree to obtain a plurality of patterns; and determining a threshold number of conditioning data smaller than the initial conservative number of conditioning data beyond which estimated facies probabilities are not significantly modified by additional number of conditioning data.

In the characterization of oil fields in the petroleum industry, 3D modeling using geostatistics is often used to assess reservoir heterogeneity and connectivity. Geostatistics often uses kriging to interpolate between data points or conditioning data. Conditioning data includes well log hard data, but can also include soft data, typically seismic data.
Conventional 3D modeling methods are based on variogram or two-point-statistics. Variogram-based algorithms allow integrating well and seismic data using a pixel-based approach. First, well data are blocked to the reservoir stratigraphic grid, i.e. well data values are assigned to the cells that the wells penetrate and sample. Then, all unsampled cells in the reservoir stratigraphic grid are simulated conditional to well and seismic data using some form of kriging. However, the models built using conventional variogram-based methods are most often not consistent with geological interpretation. Variogram-based geostatistics is inadequate in integrating geological concepts: two-point statistics variograms do not allow modeling complex geological heterogeneity. As a result, the variogram-based methods usually generate models that provide poor reservoir performance forecasting.
Over the past 10 years, the traditional variogram-based methods have been replaced by Multiple Point Statistics (MPS) methods. The MPS approach replaces traditional variograms with 3D numerical conceptual models of the subsurface geology, also known as training images.
MPS simulation is a reservoir facies modeling technique that uses conceptual geological models as 3D training images (or training cubes) to generate geologically realistic reservoir models. The training images provide a conceptual description of the subsurface geological geobodies, based on well log interpretation and general experience in reservoir architecture modeling. MPS simulation extracts multiple-point patterns from the training image and anchors the patterns to reservoir well data.
A 3D data template is provided by a user to define the dimensions of the multi-point patterns to be reproduced from the training image. Specifically, a size of the 3D data template corresponds to the maximum number of conditioning data used to infer statistics from the training image during the MPS simulation process. The larger the template (i.e., the larger the number of conditioning data), the better the reproduction of the geological features displayed by the training image. However, because computation time of MPS simulation increases exponentially with the size of the template, a trade-off is made between simulation computation time and training pattern reproduction quality. The trade-off is a function of the number of facies to model and the complexity of the training image.
One conventional approach to optimize the template size, i.e. the number of conditioning data used in MPS simulation, for a given training image is by trial-and-error where multiple MPS simulations are run using various template sizes, and the smallest template size that allows reasonable training image reproduction is retained. However, this conventional approach is extremely computation time-intensive (i.e., CPU time-intensive) and seldom used. Instead, common conventional methods consist of setting, by the experienced modeler, the template size (i.e., the maximum number of conditioning data) to an arbitrary number that is conservative and high enough to ensure reasonable pattern reproduction, even with complex multi-facies training images. For example, in some cases, the default template size is set to 40 cells. In this case, pattern reproduction is, in most cases, satisfactory. However, one drawback is that the simulation can take up to several hours (in some instances tens of hours).
Therefore, there is a need for a method for optimizing conditioning data templates that cures these and other deficiencies in the conventional methods so as to find an optimal template size, or optimal number of conditioning data, that would minimize computation time of an MPS simulation while preserving training pattern reproduction quality.

System and method for optimizing the number of conditioning data in multiple point statistics simulation


Structures for a photonics chip that include a fully-depleted silicon-on-insulator field-effect transistor and related methods. A first device region of a substrate includes a first device layer, a first portion of a second device layer, and a buried insulator layer separating the first device layer from the first portion of the second device layer. A second device region of the substrate includes a second portion of the second device layer. The first device layer, which has a thickness in a range of about 4 to about 20 nanometers, transitions in elevation to the second portion of the second device layer with a step height equal to a sum of the thicknesses of the first device layer and the buried insulator layer. A field-effect transistor includes a gate electrode on the top surface of the first device layer. An optical component includes the second portion of the second device layer.

The present invention relates to photonics chips and, more specifically, to structures for a photonics chip that include a fully-depleted silicon-on-insulator field-effect transistor and methods of fabricating a structure for a photonics chip that includes a fully-depleted silicon-on-insulator field-effect transistor.
Photonics chips are used in many applications and systems including, but not limited to, data communication systems and data computation systems. A photonics chip integrates optical components, such as waveguides, photodetectors, modulators, and optical power splitters, and electronic components, such as field-effect transistors, into a unified platform. Among other factors, layout area, cost, and operational overhead may be reduced by the monolithic integration of both types of components on the same chip.
A photonics chip may be fabricated using a silicon-on-insulator substrate as a starting platform. The silicon-on-insulator substrate includes a device layer, a handle substrate, and a buried oxide layer that electrically isolates the device layer from the handle substrate. The performance of optical components may be limited because of restrictions imposed on the thickness of the device layer by the electronic components. For example, fully-depleted field-effect transistors are not found in a photonics chip due to the large device layer thickness. In addition, the ability to implement back-gate control for field-effect transistors is limited by the ability of the buried oxide layer to provide adequate electrical isolation.
Improved structures for a photonics chip that include a fully-depleted silicon-on-insulator field-effect transistor and methods of fabricating a structure for a photonics chip that includes a fully-depleted silicon-on-insulator field-effect transistor are needed.

Photonics chips including a fully-depleted silicon-on-insulator field-effect transistor


Disclosed is an image inputting apparatus including: a base; a sensor main body coupled with the base so as to be able to move in a vertical direction on the base; a switch provided on either one of an upper surface side of the base or a lower surface side of the sensor main body; and a convex part provided on another one of the upper surface side of the base or the lower surface side of the sensor main body at a position facing to the switch.

1. Field of the Invention
The present invention relates to an image inputting apparatus, and more particularly to a switch mechanism to indicate a start of image inputting.
2. Description of Related Art
In recent years, the requirements of reliability and safety to the management and the protection of personal information or business information, or to the security thereof such as the prevention of crimes, have socially risen, and person authentication has been recognized as a very important subject. In order to meet the rise of such a social requirement, the technology developments of the person authentication using biological information, such as a fingerprint, an iris, or a vein, have been actively performed.
In order to take the biological information into a person authentication system, various sensors have been developed. As the sensors, there are optical, electrical capacitance, electric field, thermal, pressure sensitive, ultrasonic, and the other sensors. The surface or internal morphology of a testing object (for example, a finger) that touches the contact surface of a sensor is taken in as an image by the sensor, and the take image is used as biological information to be collated by the person authentication system.
In order to suppress the power consumption of the sensor, the sensing operation of the sensor is performed only when the testing object touches the sensor. To put it concretely, a push switch is provided on the lower side of the sensor, and an output signal of the push switch operates as a trigger to start the operation of the sensor when the sensor is pushed down by the testing object and the push switch is pushed (see, a first patent document: Japanese Patent Application Laid-Open Publication No. 2000-57328).
Now, the technique disclosed in the first patent document disposes a leaf spring 15 urging a sensor device 2 upward in a space between the sensor device 2 and a contact connection switch 18, and allows the leaf spring to touch the contact connection switch against the urging force thereof to conduct only when the sensor device is pushed in by the testing object. Thus, a sensing operation is performed. However, in the technique, the leaf spring is bent upward at the time of non-sensing operation (see FIG. 1 of the first patent document). Even if the leaf spring receives depressing force by the testing object at the time of a sensing operation, the leaf spring is not always surely flattened to touch the contact connection switch to conduct (see FIG. 3 of the first patent document).
Furthermore, if the sensor is further pushed in while the push switch is pushed by the sensor, then the sensor is bent in the state of being supported by the push switch. Consequently, there is the possibility of the damage of the sensor. Accordingly, the present applicant has zealously performed research and development, and suppressed the damage of the sensor by attaining the improvement of the strength of the sensor. But it is still desired to reduce the probability of occurrence of the damage of the sensor furthermore.


This invention forms a loose fill packing material comprising clean, elongated strips of paper formed from corrugated cardboard having internal fluting formatted to expand in height, improving the ability of the strips to lattice. A vertical slitter divides cardboard into blanks of a predetermined width generally with the grain of corrugation. The blanks are fed into a shredding device that cuts the blanks with a scissor motion across the width of the blank. The strips pass onto a cleated perforated conveyor and into an enclosed suction housing. Inside the suction housing is a vortex box enhancing air flow for cleaning the strips. Strips are centered on the conveyor. The strips move from the conveyor to an angled, vibrating sifter plate into a collector. Strips are sprayed with a liquid material having microbicidal, sanitizing, insect repellant, disinfectant and/or deodorizing properties.

1. Field of the Invention
This invention relates to an apparatus for converting corrugated cardboard boxes into loose fill packing material comprising elongated strips of cardboard which are substantially free of dirt, dust and other contaminants, and, which are coated with a liquid material having microbicidal, sanitizing, insect repellant, disinfectant and deodorizing properties.
2. Description of Related Art
There is a wide variety of packing materials, including Styrofoam peanuts. However, items packed in peanuts tend to shift, settle or creep within a container in the course of shipment. Also, the peanuts tend to move relative to one another, thereby permitting the packed items to shift.
This problem of shifting has been solved to some extent by the use of other packing materials in addition to or as a replacement for the Styrofoam peanuts. For example, "bubble-wrap" has also been used as an addition to or a substitute for Styrofoam peanuts. Unfortunately, bubble-wrap is also expensive and can be difficult to work with depending upon the size, shape and/or weight of an item to be shipped.
An issue with foam packing materials is that they promote the formation of static electricity within the container. The static can create substantial problems with sensitive electrical components. Additionally, electrostatically charged material attracts contaminants and other impurities which make their use undesirable for certain hygienic applications. Furthermore, the statically charged materials are difficult to handle, as the static charge causes the pieces of fill to repel each other.
Plastic packing materials create a disposal problem and can be dangerous to the environment. They use resources to create and do not readily biodegrade when placed in a landfill. Additionally, leaching can take place with certain plastics, thereby creating environmental hazards. Several states have considered laws to reduce or eliminate the dumping of plastics within sanitary landfills and this could pose a significant disposal problem. Also, plastic packing materials are non-absorbent if liquid leakage occurs during shipment. Another problem with Styrofoam is that they are impractical when shipping heavy objects, because the pieces of fill may break.
Also, cardboard strips used to protect packaged goods have several problems associated with it. However, current methods are inefficient in the manufacture of the fill. In addition, a large amount of tabs and odd shaped pieces are formed with the fill, causing interference with the latticing effect of the strips and breakdowns and increased maintenance of the apparatus. In addition, Furthermore, the dirt and dust raised by current methods make the fill material inappropriate for applications where cleanliness in the material is desired.

Loose fill packing material and apparatus for manufacturing same


A phase-locked loop circuit, in which it takes a very short time until the internal clock becomes stable while switching external clocks, is provided with a phase comparison part, a voltage controlled oscillating part, plural dummy clock producing parts, a selection supply part and a monitor control part. The phase comparison part outputs a voltage corresponding to a phase difference between two inputted clocks. The voltage controlled oscillating part outputs an internal clock of a frequency corresponding to the voltage outputted from the phase comparison part. The plural dummy clock producing parts produce dummy clocks by dividing an internal clock outputted from the voltage control oscillating part, respectively, while resetting the respective dividing operations in synchronous with inputted external clocks. The selection supply part supplies an external clock selected among plural external clocks inputted into said plural dummy clock producing parts and a dummy clock outputted from a dummy clock producing part receiving the external clock to the phase comparison part. The monitor control part monitors the states of the plural external clock and controls the selection supply part in a manner that a normal external clock is inputted into the phase comparison part. Moreover, the monitor control part controls the dummy clock producing part which the normal external clock is inputted in a manner that a dummy clock is produced without the reset of the dividing operation.

1. Field of the Invention
The present invention relates to a PLL (Phase-Locked Loop) circuit to synchronize an internal clock with an external clock from another apparatus, more particularly, relates to a PLL circuit capable of synchronizing an internal clock with one of plural external clocks for an electronic automatic exchange and the like.
2. Description of the Related Art
A PLL circuit produces a clock synchronizing with a clock from an external apparatus and including no jitter. For example, the PLL circuit is installed in an electronic automatic exchange to produce a communication system clock which synchronizes with an AMI clock (hereinafter, called external 64k+8k clock) to be a synthetic clock of a clock of a 64 kHz frequency (hereinafter, called external 64k clock) and a clock of a 8 kHz frequency (hereinafter, called external 8k clock) transmitted from a network such as the INS network.
Further, there are many cases in that a PLL circuit is fabricated to be capable of receiving clocks from plural supply sources in consideration of troubles at a network to be an external clock supply source.
FIG. 7 illustrates a sample of an above described conventional PLL circuit. As shown in FIG. 7, the conventional PLL circuit is provided with receiver circuits (RCV) 31.sub.1 -31.sub.4, an input monitor circuit 33, an input selection circuit 34, a phase comparison circuit 35, a low-pass filter (LPF) 36, a voltage controlled oscillator (VCO) 37, a 8k clock producing circuit 38 and a clock driver (CLKDRV) 39.
The RCVs 31.sub.1 -31.sub.4 separate the 64k+8K clock into the external 64k clock and the external 8k clock. The RCVs 31.sub.1 -31.sub.4 receive the 64k+8K clocks (IN1-IN4) from various clock supply sources, respectively. The external 64k clocks separated in the RCVs 31.sub.1 -31.sub.4 are supplied to the input monitor circuit 33, and the external 8k clocks are supplied to the input selection circuit 34.
The input monitor circuit 33 outputs one input selection signal indicating an IN in a normal status among four input selection signals corresponding to the IN1 through the IN4. The external 64k clocks separated in the RCVs 31.sub.1 -31.sub.4 are used in the input monitor circuit 23 for determining which input selection signal is outputted. That is, the input monitor circuit 33 grasps the status of each IN by monitoring always the status of each external 64k clock (the presence of each signal), and outputs the input selection signal indicating one IN in the normal status. Moreover, in outputting the input selection signal indicating an IN, when it is detected that this IN becomes defective, the input monitor circuit 33 specifies one IN among the other INs in normal status and outputs the input selection signal indicating the specified IN into the input selection circuit 34.
The input selection circuit 34, what is called a 4-1 selector, outputs the external 8k clock from the RCV 31 corresponding to the IN indicated by the input selection signal from the input monitor circuit 33 to the phase comparison circuit 35.
The phase comparison circuit 35, in addition to this external 8k clock, receives a 8k clock for phase comparison to be a signal from the 8k clock producing circuit 38 (details is described later). The phase comparison circuit 35 outputs a pulse signal having a width and a polarity corresponding to the phase difference between these two clocks into the LPF 36.
The LPF 36 integrates the inputted pulse signal to thereby convert the pulse signal to a voltage signal including no high frequency component. Then, the LPF 36 supplies the voltage signal to the VCO 37 as a frequency control voltage.
The VCO 37 produces a clock of a frequency corresponding to the frequency control voltage, and produces the clock of the frequency of 16.384 MHz when the frequency control voltage of "2.5" V is supplied. The clock produced in the VCO 37 (hereinafter, called 16M clock or VCO clock) is supplied to the 8k clock producing circuit 38 and the CLKDRV 39.
The 8k clock producing circuit 38 produces a 8k clock by dividing the 16M clock produced in the VCO 37, and outputs it. The output from the 8k clock producing circuit 38 is inputted into the phase comparison circuit 35 and the CLKDRV 39.
The CLKDRV 39 supplies the 8k clock from the 8k clock producing circuit 38 and the 16M clock from the VCO 37 to a system (not shown) such as a communication system.
That is, the PLL circuit shown in FIG. 7 is fabricated to adjust the frequency of the VCO clock in a manner that the phase difference between the 8k clock outputted from the 8k clock producing circuit 38 and the external 8k clock supplied to the phase comparison circuit 35 becomes "0".
This PLL circuit produces an 8k clock including no jitter and coinciding with the external 8k clock concerning a phase and a frequency, therefore, a satisfied communication can be obtained. Moreover, when a trouble occurs in the IN including that external 8k clock, the IN is switched automatically and an 8k clock synchronized with the switched IN is produced, so that the communication can be continued.
In the conventional PLL circuit, however, there are some case that a trouble occurs in a system or a communication between systems when the IN is switched.
That is, when the external 8k clock inputted into the phase comparison circuit 35 is switched as the result which a trouble occurs in an IN, usually the phase of the external 8k clock before switching (i.e. the 8k clock for phase comparison outputted from the 8k clock producing circuit at that time) does not coincide with the phase of the external 8k clock after switching, therefore, the frequency of the VCO clock is adjusted in a manner that the phase difference between the 8k clock outputted from the 8k clock producing circuit 38 and the external 8k clock after switching becomes "0".
For instance, there may be a maximum phase difference of 62.5 .mu.s between the two external 8k clocks. Thus, in the worst case, immediately after changing over the clocks, there is the phase difference of 62.5 .mu.s between the external 8k clock and the 8k clock for phase comparison inputted into the phase comparison circuit 35. Assuming that the frequency variation range of the VC0 37 in the PLL circuit shown in FIG. 7 is .+-.50 ppm of the oscillating frequency, namely, about .+-.19 Hz, it takes a time to establish the synchronization after occurring this worst case, as described later.
In this case, since the phase comparison circuit 35 outputs the maximum (or minimum) frequency control voltage, the VCO 37 outputs the VCO clock of 16.384819 MHz which is the maximum oscillating frequency. The time for one clock of this VCO clock is "1/16.384819" MHz.sup.-1 =.+-.sec!, and the time for one clock of the VC0 clock at the central frequency is "1/16.384" .+-.sec!. Thus, the phase difference between the external 8k clock and the 8k clock for phase comparison becomes smaller 3.054586.times.10.sup.-12 sec! (1/16.384-1/16.384819 .+-.sec!) for one clock of the VCO clock. Therefore, it is necessary for the VC0 clock of at least 20461038 (=62.5.times.10.sup.-6 sec!/3.054586.times.10.sup.-12 sec!) until a synchronization is established. That is, it takes at least 1.24878 (=1/16.384819MHz!.times.20461038)sec! to establish the synchronization.
A PLL circuit is also known, in which the external 64k clock is inputted into the input selection circuit 33 and the external 8k clock is inputted into the input monitor circuit 23, namely, a PLL circuit which a synchronization is established with the 64k clock. When the 64k clock is used, the maximum phase difference becomes 7.8 .mu.sec. Thus, when such a PLL circuit is fabricated with the above VCO, it takes at least 0.15609 sec! to establish the synchronization in the worst case.
As above described, in the conventional PLL circuit, when the IN is switched, it takes a time corresponding to the phase difference between the external clocks before and after switching to establish the synchronization between the external clock and the clock for the phase comparison. Therefore, in the conventional PLL circuit, there are some cases that unstable clocks are supplied to a system for a relative long time during clock switching, as the result, there are some cases that a trouble occurs in a system or a communication between systems.


An arrangement for mounting swing cylinders in a boom for a rock drilling unit which includes a frame, a boom pivotally connected about vertical and horizontal shafts, relative to the frame, a lift cylinder connected between the frame and the boom, and a swing cylinder connected between the frame and the boom. The unit also includes a support for a feed beam, pivotally connected to the other end of the boom about horizontal and vertical shafts, a tilt cylinder, connected between the support and the boom, and a transverse swing cylinder connected between the support and the boom. The swing cylinder is coupled at an angle (.alpha.), relative to the longitudinal axis of the boom, such that the longitudinal axis of the swing cylinder has a downward inclination from the frame toward the opposite end of the boom, so that with the boom in the middle of its upper and lower angles of altitude, the swing cylinder is essentially horizontal.

The invention relates to an arrangement for mounting swing cylinders in a boom for a rock drilling unit, comprising a frame, a boom pivotally connected about vertical and horizontal shafts, relative to the frame, a lift cylinder connected between the frame and the boom for vertical lifting and lowering of the boom, a swing cylinder connected between the frame and the boom for lateral turning of the boom relative to the frame, a support for a feed beam, connected to the other end of the boom pivotally about horizontal and vertical shafts, a tilt cylinder, connected between the support and the boom for turning the support relative to the boom about a horizontal shaft, and a transverse swing cylinder connected between the boom and the support for turning the support relative to the boom about a vertical shaft.
A problem with booms for rock drilling units is that when a feed beam is vertically turned when simultaneously being laterally deviated from its upright position to either side, the end of the boom moves so as to describe an arc. The greater the upward or downward movement of the boom, the greater the simultaneous lateral outward swing. This is caused by the articulated structures of a boom and the consequent structural geometry, the elimination of which has in practice proved impossible by mechanical means. A weakness of known solutions is also that the upward and downward swing angles of booms are not symmetrical; the downward angle is narrower than the upward angle for practical reasons. Consequently the lateral deviation in both angular positions of a boom is far too large to be handled and causes much inconvenience when the boom is used. When the boom is allowed to turn to its extreme lateral angles being in its upper angular position, the joints of the cylinders may shift to a position where the boom no longer can turn back, but is locked in place.
The object of this invention is to provide an arrangement with optimum elimination of the above problems and with as equal a lateral swing as possible in the upper and lower vertical positions of a boom. The arrangement of the invention is characterized in that a swing cylinder between a frame and a boom is coupled at an angle relative to the longitudinal axis of the boom so that the longitudinal axis of the swing cylinder has a vertical downward inclination from the frame toward the end of the boom, relative to the direction of the longitudinal axis of the boom, and that the centres of the horizontal joints, relative to the base, of the swing cylinder and the boom are essentially united with each other at the base of the boom when the boom faces straight ahead relative to the frame.
An essential idea of the invention is that the swing cylinder is coupled vertically inclined relative to the longitudinal axis of the boom so that with the swing cylinder in the horizontal plane, the boom is in the middle of its vertical swing range, whereat the upward or downward swinging of the boom from this position causes an equal deviation in both the extreme upper and lower positions to the lateral swing of the boom because of the position of the swing cylinder. Correspondingly, the transverse swing cylinder needed for turning the feed beam is coupled at an angle to the longitudinal axis of the boom, whereat the turning of both the feed beam and the boom results in equal angular changes by the action of the vertical swing angle.
It is an advantage of the invention that the changes in the direction and lateral deviations of both the boom and the feed beam, when vertically turned, are essentially the same relative to the horizontal plane both above and below, as the angles between the longitudinal axis of the boom and the swing cylinder, and correspondingly the transverse swing cylinder, compensate for the vertical deviations of the boom.
The invention will be described in greater detail below, in connection with the accompanying drawings.

Arrangement for mounting swing cylinders in boom for rock drilling unit


Coffee liquid extracting apparatus comprises an extractor vessel, a combined-use decanter that retains water for extraction, receives extracted coffee liquid and is usable for pouring coffee liquid into cups. A circulation path and a circulation pump transport the liquid from the combined-use decanter and pour the liquid into the extractor vessel. A circulating water heater heats water circulated by the circulation pump. A controller controls operation in response to a tap temperature operation switch to operate the circulation pump to circulate the water for a fixed time (T.sub.1) with the circulating water heater on for a fixed time (T.sub.0) limited to an early operation period, thereby pouring high temperature water for a fixed time (T.sub.0) limited to the early operation period and subsequently pouring unheated water.

The present invention relates to a coffee liquid extracting apparatus.
There are two methods for obtaining ice coffee. One method is to first extract hot coffee liquid using high temperature water and then to cool the extracted coffee liquid. Another method is to extract coffee liquid by the use of ambient temperature water, such as unheated, or ordinary tap temperature water, from the beginning, and then to cool the extracted coffee liquid to a low temperature.
However, in the above former method, wherein coffee liquid is extracted by the use of high temperature water, cooling to a low temperature requires a large amount of cooling and thus consumes a lot of ice in order to change the coffee liquid extracted by high temperature water into ice coffee. In actual practice, such a large amount of cooling requires ice making machinery.
As opposed to and as compared with the former method, the latter tap temperature water method can reduce an extracted quantity of such components as astringency, caffeine, sundry taste and the like that deteriorate the taste of coffee liquid. Further, the latter method includes the merit of decreasing oil content which causes oxidation of coffee liquid. However, the problem of the latter method is that a long time for extraction is required. For instance, in existing methods, about one night, namely at least 5 to 8 hours, is required to extract coffee liquid by the use of tap temperature water.

Coffee making apparatus for making coffee using heated and unheated water


A tomographic imaging system, including a data processing component having a memory and at least one processor configured to:                access scattering parameter data representing electromagnetic waves scattered by features within an object and originating from a plurality of antennas disposed around the object on a boundary S;        process the scattering parameter data to generate a reconstructed image representing a spatial distribution of features within the object, said processing including:        solving an electromagnetic inverse problem, wherein forward and inverse steps of the inverse problem are represented and solved as respective differential equations involving an electric field to determine values for the electric field; and        process the determined values of the electric field to generate reconstructed image data representing one or more spatial distributions of one or more electromagnetic properties within the object.

Inverse problems arise in electromagnetics (EM) whenever the effects of a cause are known, but the cause itself is unknown. The aim of an inverse electromagnetic problem is thus to find the EM properties of the cause according to the observed effects. By contrast, in forward problems the cause is known and the desire is to calculate the corresponding effects. In a wide range of electromagnetic inverse problems, e.g. geophysical or biomedical imaging, the cause is usually the complex permittivity (containing the relative permittivity and the conductivity) which describes an object electromagnetically, while the effects are the scattered electric fields from the object, which are detected by imaging antennas and represented by what is referred to in the art as "scattering parameters" or "S-parameters", as described in W. C. Chew, Waves and Fields in Inhomogeneous Media, IEEE Press, New York, 1995 ("Chew").
The main difference between forward and inverse problems is that a forward problem has a unique solution, whereas an inverse problem does not have a unique solution because the available data is insufficient to determine a unique solution.
For example, when a dipole antenna (cause) with a certain port impedance, given structure and substrate, and specific operating frequency radiates electromagnetic waves in free space, there is only one possible field distribution (effect) formed in free space by this antenna. However, if this field distribution (effect) is known and the desire is to find the properties of the antenna (cause) that has generated this field distribution, in theory an infinite number of corresponding antennas can be found.
The reconstructed images of quantitative EM inverse problems (tomography) are usually coarse and inaccurate due to diffraction effects, which are significant at microwave frequencies (unlike high frequencies such as X-Rays which propagate in straight lines). Nevertheless, these coarse images are useful as they are reconstructed from a domain non-invasively. For example, in reconstructing an image from the inside of a 'black box' (e.g., a human head, an underground mine, etc), having a coarse image representing the spatial distribution of complex permittivity within the black box is much more important than a requirement for high quantitative accuracy. This point is even more important in biomedical cases where conventional techniques such as MRI or CT-Scan are undesirably time-consuming, static and bulky, or have ionization radiation, are not portable and thus cannot be used on-site in emergency scenarios such as traffic or sport accidents or stroke where time generally has an inverse relation with the chance of survival.
In inverse problems, the permittivity over the domain of interest is always described by three-dimensional (3D) integral equations. However, as described in Chew, these equations are too complex to be solved directly, and therefore are always approximated by two-dimensional (2D) integrals to make them solvable. A two-dimensional image representing the dielectric properties of the domain is reconstructed through an iterative forward-inverse computational procedure using either Born's method, Van Den Berg's (different contrast source inversions) method, or gradient-based methods.
It is well known that these current state of the art methods for solving electromagnetic inverse problems suffer from serious limitations, including the requirements to represent the imaging antennas as ideal point sources, to assume homogeneity along one of the coordinate axes, and to use a background matching medium. These limitations restrict the reconstruction of images to two dimensions with a certain level of accuracy error, and impose manufacturing problems associated with selecting a suitable background medium with reasonable dielectric properties. Additionally, to reduce antenna mutual coupling and surface waves, a lossy matching medium is usually used. However, this also attenuates the desired signal by more than 10 dB (i.e., 90% of the useful signal will be lost), and thus the dynamic range and detection capability are significantly limited. Additionally, these methods can also be undesirably time-consuming, especially for biomedical imaging applications in emergency scenarios, such as stroke or trauma patients, where the image reconstruction time can have a dramatic effect on patient outcomes, and might even mean the difference between life and death in some instances.
It is desired, therefore, to overcome or alleviate one or more difficulties of the prior art, or to at least provide a useful alternative.


A surgical instrument for removing a cataract from a patient's eye including a hollow vibrating needle surrounded by two hollow infusion sleeves. The outer sleeve conforms to the incision in the eye and the inner sleeve prevents the outer sleeve from collapsing onto the hollow, vibrating needle.

This invention relates generally to a surgical instrument for use in eye surgery, and, more particularly, relates to a surgical instrument for cataract eye surgery.
A wide array of fluid-irrigated, ultrasonically-operated cutting devices have been developed for ophthalmological surgical techniques such as phacoemulsification--a method for removing a cataract through a surgical incision in the eye. Examples of such devices are set forth in the following patents:
______________________________________ U.S. Pat. No. Inventor ______________________________________ 3,589,363 Banko et al. 4,515,583 Sorich 4,808,154 Freeman 4,816,017 Hood et al. 4,816,018 Parisi ______________________________________
U.S. Pat. No. 3,589,363 to Banko et al. and U.S. Pat. No. 4,515,583 to Sorich disclose ultrasonic surgical instruments that utilize rigid outer sleeves.
U.S. Pat. No. 4,808,154 to Freeman discloses a phacoemulsification/irrigation and aspiration sleeve apparatus.
U.S. Pat. No. 4,816,017 to Hood et al. discloses an ultrasonic decoupling sleeve.
U.S. Pat. No. 4,816,018 to Parisi discloses an ultrasonic probe tip.
Fluid-irrigated, ultrasonically-operated cutting devices typical of the prior art, however, suffer from a number of deficiencies. In particular, most infusion sleeves used during phacoemulsification are made of silicone or a silicone-type material. The use of this type of infusion sleeve can cause fluid leakage between the incision edge in the eye and the exterior surface of the infusion sleeve, resulting from a need to make the incision in the eye larger than the infusion sleeve. This need is due to the compressibility of silicone or like materials which cannot be safely used when inserted through an incision in the eye where there is a minimal amount of clearance between the incision and the exterior of the silicone infusion sleeve.
When there is a minimal clearance between the exterior of the silicone infusion sleeve and the incision of the eye, the incision tends to compress the non-rigid silicone infusion sleeve against the vibrating tip which results in relative rubbing movement between the silicone sleeve and the vibrating tip. This relative movement generates undesirable heat as the needle is being vibrated at its relative high frequencies. The generation of this heat is extremely undesirable inasmuch as it can result in thermal burns and shrinkage of the ocular tissue surrounding the silicone infusion sleeve. The burning and shrinkage of ocular tissue is a serious problem with sight-threatening implications.
In an attempt to reduce the infusion fluid leakage and the deleterious effects that can be caused by undesirable friction generated therefrom, some infusion sleeves have been constructed from rigid, non-compressible materials. Generally these materials have consisted of Teflon or metallic-based compositions. These rigid, infusion sleeves have been relatively successful in solving the problems of constriction of the path for fluid flow between the distal end of the infusion sleeve and the vibrating tip as well as the heat generation and thermal burns associated with malleable infusion sleeves; however, other problems persist with these non-compressible infusion sleeves.
While rigid sleeves are capable of being inserted through smaller incisions, which has the advantage of reducing leakage, there is still significant leakage. The primary cause of the persistent leakage between the rigid infusion sleeve and the eye incision is that the cross section of the rigid sleeve does not match the contour of the eye incision. As a consequence, there are fairly substantial gaps between the rigid sleeve exterior surface and the eye incision. This is because the collagen fiber structure of the cornea resists deformity and thus does not readily assume the shape of the infusion sleeve.
The experience of the applicant, who has performed literally thousands of cataract eye operations, has shown that it is impossible, from a practical standpoint, to fully eliminate the problem of leakage during cataract surgery by means of a smaller incision and forcing the rigid infusion sleeve through it. While this may decrease wound leakage, it does not eliminate the problem and it causes the instrument to be so tightly held by the deformed incision that there is great difficulty in advancing and withdrawing the instrument through the incision. As will be apparent to those skilled in the art, during cataract surgery the instrument must be advanced and withdrawn many times through the incision as the fractured portions of the cataract are removed from the various locations within the anterior and posterior chambers of the eye.
Although there are means to prevent leakage, these means are very expensive to construct. They require precision tools with very precise configurations. The invention disclosed in U. S. Pat. No. 5,084,009 to Mackool, addresses the aforementioned problems; however, it also requires concentricity between the rigid infusion sleeve and the vibrating needle.
Accordingly, there exists a need for an improved apparatus for performing cataract surgery.
It is accordingly an object of the invention to provide an improved apparatus for performing cataract surgery.
It is another object of the invention to provide such apparatus that is relatively inexpensive to construct.
Other general and specific objects of the invention will in part be obvious and will in part appear hereinafter.


A string cutter head for weed trimmers comprises a solid seamless hub having an upper disk and a lower disk. The lower disk comprises at least two pairs of vertically drilled holes for installation of a monofilament cutting string. The upper disk acts as a spacer from the trimmer housing to allow installation and removal of cutting string without removal of the cutting head from the weed trimmer. Several holes of different diameters may be used for different diameter strings. The hub comprises a threaded connection for attachment to the shaft end of the trimmer. In the preferred embodiment, the hub is machined from aluminum alloy.

The present invention relates to weed trimmers and, more particularly, to weed trimmer cutting heads.
Weed trimmers have gained popularity as a convenient device for cutting grass, weeds, and brush, especially in difficult access areas around home, businesses, and farms. A number of cutting heads for these devices have been disclosed to meet the needs of users. An example of such a device is disclosed in U.S. Pat. No. 5,640,836 by the applicant. The blade assembly for weed trimmers disclosed in this patent comprises pivoting steel blades for cutting heavy weeds, vines and underbrush as well as grass and lighter growth.
In many residential areas, string type weed trimmer heads utilizing a monofilament plastic string are popular due to the fact that the string cutting elements are less likely to cause damage to trees, shrubs, and fences if the cutting head is brought too near the trimmed object. These devices are often simpler and less costly than blade units.
One of the major disadvantages of the string type cutting heads is that the string is easily cut or damaged, especially if hard objects are inadvertently encountered. In these cases, the string must be frequently replaced, taking considerable time. In an effort to address this problem, a number of cutting heads have been introduced on the market which contain a spool of string. A number of methods are used to manually, or in some cases automatically, extend additional string when the old string becomes cut or damaged.
Theses units, while partially successful in addressing the problem of string replacement, have added bulk and weight, making use of the trimmer more difficult. Many smaller weed trimmers do not have the power or ruggedness to utilize these units. In many cases, these units fail to function as intended due to poor design, inadequate maintenance, or improper use. Changing string types or diameters is difficult and time consuming. Inventories of different string types is expensive due to the quantities of string used in the spools. These devices also are expensive to purchase and maintain.
So called "cut length" string trimmers, on the other hand, are simple, light and reliable to use. However, many require too much time for string replacement, especially when the head must be removed from the trimmer to replace the string.


A linear magnetic damper includes a magnetically attractive rail, a damper body movable along the magnetically attractive rail and having a longitudinal passage, which accommodates the magnetically attractive rail, and having a plurality of openings symmetrically formed at two sides thereof and respectively arranged in alignment, and a plurality of magnetic members respectively mounted in the openings at two sides of the damper body for applying a magnetic force to the magnetically attractive rail.

1. Field of the Invention
The present invention relates to a damper for use in an exercising apparatus and more particularly, to a linear magnetic damper.
2. Description of the Related Art
FIG. 1 illustrates a conventional exerciser, which is comprised of a framework 2, a lifting mechanism 3 mounted on the framework 2, a plurality of pulleys 5, a weight set 6, and a cord member 4 running on the pulleys 5 and connected between the lifting mechanism 3 and the weight set 6. When lifting the lifting mechanism 3, the weight set 6 generates a counter force (resistance) against the user. Because the weight set 6 includes a plurality of heavy metal members, these have metal members may hit against one another to produce much noise during operation of the exerciser (more particularly during down stroke of the weight set 6, thereby annoying the user. It is also dangerous when the weight set 6 falls to generate great impact. Further, because the lifting mechanism is operated to fight against the gravity of the weight set, the weight set must be arranged to move along a vertical path, i.e., there is a limitation to the installation orientation of the weight set.


A flow damper for cleaning stations, in particular for cleaning pipettor tips, is described, a cleaning liquid being introduced into the cleaning station through a pipettor needle and flowing around the pipettor tip before it is conducted away out of the cleaning station. In order to avoid wetting of the isolation surfaces of the pipettor needle and a level detector electrode, if provided, there are provided, according to the invention, in the cleaning station at least one screen having a preferably central entry orifice for the jet of liquid emerging from the pipettor needle and a device for diverting the liquid flow in the direction essentially opposite that of the entering jet of liquid. Furthermore, a plurality of exit orifices, preferably arranged radially around the entry orifice, for the diverted liquid flow are formed in the at least one screen.

1. Field of the Invention
The invention relates to a flow damper for a cleaning station, in particular for cleaning pipette tips, a cleaning liquid being introduced into the cleaning station through a pipette needle, and flowing around the pipette tip before it is conducted away out of the cleaning station, and relates to a cleaning station having such a flow damper.
2. Description of the Related Art
After processing reagents or patient samples, the pipette tips which have come into contact with them must be cleaned before they can be used for processing further reagents or samples. The pipette tips are usually cleaned in overflow cleaning stations with cylindrical overflow vessels in which a cleaning liquid is injected at a high speed through the pipette needle. The cleaning fluid accumulates there and flows around and cleans the pipette needle from the outside before it is conducted away via an overflow. Owing to the high flow speeds required for the internal cleaning of the pipette needle, severe turbulences occur in the overflow vessels, generating waves, bubbles and splashes onto the surface of the vessel. As a result, the upper isolation surfaces of the pipette needle which do not have to be cleaned and of a level detector electrode used to determine the position of the pipette needle in the cleaning station, also used at other processing stations, are also covered with liquid. Apart from the fact that this covering with liquid may lead to a carry-over of reagents or patient samples, the liquid adhering to the isolation surfaces of the pipette needle and the level detector electrode adversely affects the level detection based on a potential measurement, in that it generates a short circuit which makes reliable position measurement impossible.


A variable-diameter reinforcing cage for an anchor rod or pile foundation, comprising an axial rod (4), a plurality of vertical bars (2), at least two circular fixators (5) and several groups of ribs (3) which correspond to the circular fixators (5). The circular fixators (5) are all sleeved on the axial rod (4); one end of the ribs (3) is movably connected to the position of the vertical bars (2) at the same height, and the other end of the ribs (3) is movably connected to the circular fixators (5); annular hoops (6) are arranged on the periphery of the vertical bars (2) to serve as circles of latitude; fixing points are formed on the annular hoops (6) and the vertical bars (2); and the annular hoops (6) are annular spiral spring hoops of an elastic material or flexible steel wires.

The anchor rod must have several elements: a rod body with a tensile strength higher than that of the soil-rock mass: one end of the rod body can be in close contact with the soil-rock mass to form friction (or bonding) resistance; the anchor rod body is located at the other end of the soil-rock mass and can form a radial resistance to the soil-rock mass; as a tensile member deeply penetrating into strata, the anchor rod is connected to an engineering structure at one end and deeply penetrates into strata at the other end; the whole anchor rod is divided into a free section and an anchoring section; the free section refers to the area where the tension at the anchor rod head is transmitted to the anchoring body, and the function of the free section is to prestress the anchor rod; the anchoring section refers to the area where cement slurry or concrete anchoring bodies bond prestressed reinforcements with soil layers, and the function of the anchoring section is to increase the bonding friction between the anchoring body and soil layers as well as the bearing capacity or tensile force of the anchoring body and transmit the tension of the free section deeply to the soil mass.
On the whole, the anchor rod is a member system structure for soil-rock mass reinforcing. Depending upon the longitudinal tension of the anchor rod body, the shortcoming that the tensile strength of soil-rock mass is far lower than its compressive capacity is overcome. From the viewpoint of mechanics, the cohesion C and the internal friction angle φ of the surrounding rock mass are mainly increased. In essence, the anchor rod is located in the soil-rock mass and forms a new complex with the soil-rock mass. The anchor rod in this complex is the key to solving the problem about the low tensile strength of the surrounding rock mass. Thereby, the bearing capacity of the soil-rock mass is greatly enhanced.
Anchor rods are the most basic component of the roadway support in underground mining in modern times. Anchor rods tie the surrounding rocks of a roadway together, so that the surrounding rocks support themselves. Now anchor rods are used in mines and also in construction engineering technologies, and actively reinforce basements, slopes, tunnels, dams, etc.
The basic type of anchor rods: steel bar or wire rope mortar anchor rods. Cement mortar is used as the binder for anchor rods and surrounding rocks. There are also inverted wedge metal anchor rods, pipe joint anchor rods, and resin anchor rods. The use of resin as a binder for anchor rods is costly.
Hui Xingtian of Xi'an University of Science and Technology invented a new type of spiral anchor rod→spinning anchor rod. There are multiple types of spinning anchor rods as follows: self-tapping extrusion screwing-in anchor rod→direct extrusion screwing-in and installation in soil layers without need for boring; anchoring force 20 KN/m; spinning grouting anchor rod→a spinning grouting anchor rod with initial anchoring force is formed by use of spinning anchor rod grouting after installation in a borehole; spinning resin anchor rod→a spinning resin anchor rod with initial anchoring force is formed through resin mixing at the same of installation in a borehole; self-drilling and self-anchoring rod→a self-drilling anchor rod with initial anchoring force is formed by placing drill pipe in the hollow of the spinning anchor and completing hole drilling and installation at a time; spinning shotcrete anchor rod→completion of installation, anchoring and grouting in soil layers at a time through guniting while drilling; anchoring force 35 KN/m;
At present, the expanded anchor rod technologies commonly used in the market include plain slurry expanded anchor rod technology, bag type expanded anchor rod technology, etc. In terms of construction cost, there is a foundation for the hole expanding technology for a reducer anchor rod or pile foundation. Grouting or concrete injection forms a reducer, but it does not involve a suitable reinforcement cage, a tension or resistance transmission anchor rod or pile foundation with sufficient friction cannot be formed, and especially the anchoring force of the anchor rod is limited. When this technology is used in anti-floating foundation pit supporting, slope supporting, reinforcement, etc. for building basements, the anchoring force is insufficient. This is because there is a need for large uplift resistance and high stability and reliability.
In addition, non-reducing foundation piles are often used in the foundations of high-rise buildings. On the premise of meeting the same strength and deformation requirements, in comparison with a non-reducing foundation pile, a variable-diameter foundation pile mainly has the following features: 1) at the same length, the strength of a variable-diameter foundation pile can be increased by 1.1 to 1.5 times, and its deformation can be reduced by 0.7 to 0.9 times.
2) Under the first feature condition, the strength and deformation requirements of buildings are met, and the pile length can be obviously shortened.
3) On the premise of meeting the strength and deformation requirements of the pile, the pile length is shortened, the workload is reduced, and the construction conditions are improved so as to save labor, materials and time.
The strength of some clay layers, soft layers, pebble layers, gravel layers or weathered rock formations tends to be lower than that of concrete, thus making against the bearing capacity of concrete foundation piles. Therefore, in order to give full play to the strength characteristics of concrete, there have been papers in the field of construction engineering proposing that the use of variable-diameter foundation piles to increasing the bearing capacity of piles is technically reasonable and feasible.
The calculation method for the strength and deformation of variable-diameter foundation piles is the same as that for the strength and deformation of non-reducing foundation piles. Foundation piles are divided into end bearing piles and friction piles. The calculation method for the strength and deformation of the two piles is different. Friction piles are used as the object for calculation and comparison. The strength of a friction pile generally consists of its side friction resistance and the strength of the pile end bearing layer. For a friction pile, the pile perimeter friction resistance is primary. However, for most foundation piles and especially large diameter piles, the pile end is supported by bedrocks, and the pile end bearing capacity is primary. Therefore, the calculation of the ultimate strength of the pile end bearing layer is very important, because the values obtained by different methods for calculating the strength of the pile end bearing layer vary greatly. The strength of the foundation pile bearing layer is related to the properties of rocks and soils as well as the burial depth and size of the pile foundation. It can be seen that the variable-diameter foundation pile is very promising in application, but how to obtain a feasible variable-diameter foundation pile is a problem to be solved.

Variable-diameter reinforcing cage for anchor rod or pile foundation

JIANGSU GREEN RIVER ENVIRONMENTAL TECHNOLOGY CO., LTD.


Methods and apparatus for enhanced overlay state maintenance in a peer-to-peer overlay network. A first method includes inferring that a first node is leaving the overlay network, and transmitting a decrement message to decrement a size counter value. A second method includes identifying a set of nodes associated with a first node of an overlay network, obtaining a segment length associated with each node of the set of nodes, and determining a size of the overlay network by dividing the total number of nodes in the set of nodes by the sum of the segment lengths. A third method includes identifying a set of nodes associated with a first node of an overlay network, obtaining a size estimate associated with the first node and with each node of the set of nodes, and determining a size of the overlay network by averaging the size estimates.

1. Field
The present application relates generally to the operation of overlay networks, and more particularly, to methods and apparatus for enhanced overlay state maintenance in a peer-to-peer overlay network.
2. Background
A network in which member nodes obtain services in the absence of server-based infrastructure is referred to as a "peer-to-peer" overlay network. In a peer-to-peer overlay, peer nodes co-operate with each other to provide services and to maintain the network. Peer-to-peer overlay networks can be built on top of an underlying network, such as a network utilizing the Internet Protocol (IP).
In a peer-to-peer overlay network, each node has knowledge of one or more peers participating in the overlay. A simple but inefficient approach to routing data on the overlay network from a source node to a destination node is to continually pass the data to the next hop or successor node (i.e., node logically next to the source node in the identity space) until the destination node is reached. However this approach incurs excessive latency as the size of the overlay network and number of hops increases. Therefore, for routing optimization, each node maintains a list of fingers that are one node, two nodes, four nodes, or up to 2(m-1) nodes away from itself; where m is the number of bits assigned to each node identifier. This helps in minimizing the number of hops needed to route data from node O(n) to node O(log(n)) where n is the number of nodes for the average case.
However, nodes may come and go at any time and result in changes in the overlay network configuration that affects the fingers known to a particular node. To compensate for this variability, each node re-runs a finger stabilization algorithm that re-computes the fingers known to that node. The results of the finger stabilization algorithm are stored in a finger table. Unfortunately, since each node maintains it own finger table, determining the frequency at which each node needs to run its finger stabilization algorithm may be problematic. For example, if a node does not run its finger stabilization algorithm often enough, its finger table may become stale resulting in inefficient and delayed packet routing. If a node runs its finger stabilization algorithm too often, this may result in wasting overlay bandwidth and/or placing a burden on other nodes on the overlay network. For example, running the stabilization algorithm requires power, and excessive execution of the stabilization algorithm may waste power at battery operated nodes.
Furthermore, since structured peer-to-peer overlay networks are highly distributed in nature, participating nodes do not have complete routing tables and hence do not know the size of the overlay network they are operating in. However, knowledge of the size of the overlay network can be useful for several purposes, such as when merging overlays, performing load balancing or caching strategies, as well as routing protocol specialization.
Unfortunately, conventional systems may fail to provide an accurate overlay network size. For example, nodes may leave the overlay without gracefully notifying their neighbors, and as a result, the size of the overlay network may not be accurately maintained.
Therefore, it would be desirable to have a simple cost effective mechanism that operates to allow a node to adaptively schedule a finger stabilization algorithm and to determine the size of an overlay network to overcome the problems described above.


This invention relates to a personal computer having facility for ready variation in the amount of memory capability available on the microprocessor local bus. The personal computer has a microprocessor, a high speed local bus coupled to said microprocessor, a system bus, a bus controller coupled to and providing communication between the local bus and system bus, and a local bus memory facility coupled to and physically separable from the local bus. The local bus memory facility is provided by a substrate, preferably a printed circuit card, for mounting and providing coupling connections among components mounted thereon, at least one volatile memory component mounted on the card, and a memory controller mounted on the card and coupled to the volatile memory component, local bus and bus controller, with the memory controller regulating communications between the volatile memory component and microprocessor.

Personal computer systems in general and IBM personal computers in particular have attained widespread use for providing computer power to many segments of today's modern society. Personal computer systems can usually be defined as a desk top, floor standing, or portable microcomputer that consists of a system unit having a principal system processor and associated volatile and non-volatile memory, a display monitor, a keyboard, one or more diskette drives, a fixed disk storage, and an optional printer. One of the distinguishing characteristics of these systems is the use of a motherboard or system planar to electrically connect these components together. These systems are designed primarily to give independent computing power to a single user and are inexpensively priced for purchase by individuals or small businesses. Examples of such personal computer systems are IBM's PERSONAL COMPUTER AT and IBM's PERSONAL SYSTEM/2 Models 25, 30, 50, 60, 70, 80, 90 and 95.
These systems can be classified into two general families. The first family, usually referred to as Family I models, use a bus architecture exemplified by the IBM PERSONAL COMPUTER AT and other "IBM compatible" machines. The second family, referred to as Family II Models, use IBM's MICRO CHANNEL bus architecture exemplified by IBM's PERSONAL SYSTEM/2 Models 50 through 95. The Family I models typically have use the popular INTEL 8088, 8086 or 80286 microprocessor as the principal processor. These processors have the ability to address one megabyte of memory. The Family II models typically use the high speed INTEL 80286, 80386, and 80486 microprocessors which can operate in a real mode to emulate the slower speed INTEL 8086 microprocessor or a protected mode which extends the addressing range from 1 megabyte to 4 Gigabytes for some models. In essence, the real mode feature of the 80286, 80386, and 80486 processors provide hardware compatibility with software written for the 8086 and 8088 microprocessors.
More recently, the technology embodied in the Intel microprocessors used in personal computers has been proposed for use in more advanced computer environments, such as tightly or loosely coupled parallel or multiple processing arrays. As the higher speed microprocessor technology is applied to such arrays, problems heretofore encountered in personal computer technology as here defined may be anticipated as appearing in such environments.
With the advent of newer operating systems designed in part to use the greater memory addressing capabilities of higher speed microprocessors there has come a need to provide facility for expansion of memory into the higher ranges. Such expansion is being accomplished using devices known as single inline memory modules, or SIMMs, made up from memory modules having varying memory capabilities. While such expansion has been accomplished (in part using strategies described more fully hereinafter), the use of ever faster microprocessors causes quality and time-of-flight of signals exchanged in memory access to become more critical. By "time-of-flight" is meant the amount of time necessary for a signal to be exchanged, such as for a data or address or other control signal to be exchanged between physical memory modules and the principal processor.
Looking toward the expansion of memory which may become desirable as an existing personal computer is adapted to developing technology, problems arise as SIMMs are used in accommodate such needs. In particular, while SIMMS of varying capacity such as 1, 2, 4 or 8 megabytes and using an assortment of different memory modules may be inserted onto a local memory bus, the number of sockets required to physically hold such SIMMs and variations among them lead to impedance mismatching. Solutions for such impedance mismatching are available, and require increased buffering (leading to slowing of memory accessing or time-of-flight) or more advanced semiconductor technology (leading to increased expense).
An alternative to SIMMs used in some computing environments is the SIPP (Single In-line Pin Package) memory module. The difference between a SIMM and a SIPP arises from the form of connector used--SIMMs use card edge connectors, while SIPPs use header pins. While the two forms are not interchangeable in terms of installation into a computer system environment, similar types of problems arise in systems using memory expansion with SIMMs and those using SIPPs.
An alternative approach is the use of memory expansion cards or boards accessed through the portion of the bus architecture provided for installation of option cards or boards such as is known for example, in the more advanced personal computers described above, as the MICROCHANNEL bus. While such an approach is viable, the installation of additional memory capability by installation of an option card uses up a scarce resource (slots on the bus) and slows access to the added memory. Further, such additions require more expensive data drivers which potentially increase electromagnetic compatibility emissions and which, in turn, lead to the use of preventative measures which cause signal under- or over-shoot and timing skew problems.

System having microprocessor local memory expansion capability through removable coupling of local memory expansion boards directly to the high speed microprocessor local bus


Crystalline solid forms of the selective late sodium current inhibitor 4-(pyrimidin-2-ylmethyl)-7-(4-(trifluoromethoxy)phenyl)-3,4-dihydrobenzo[f][1,4]oxazepin-5(2H)-one (Compound I) were prepared and characterized in the solid state:

The present disclosure relates generally to crystalline solid forms of the compound 4-(pyrimidin-2-ylmethyl)-7-(4-(trifluoromethoxy)phenyl)-3,4-dihydrobenzo[f][1,4]oxazepin-5(2H)-one, processes for making the forms, and their therapeutic methods of use.
The late sodium current (INaL) is a sustained component of the fast Na+ current of cardiac myocytes and neurons. Many common neurological and cardiac conditions are associated with abnormal INaL enhancement, which contributes to the pathogenesis of both electrical and contactile dysfunction in mammals. See, for example, Pathophysiology and Pharmacology of the Cardiac "Late Sodium Current", Pharmacology and Therapeutics 119 (2008) 326-339. Accordingly, compounds that selectively inhibit INaL in mammals may therefore be useful in treating such disease states. Such conditions include, but are not limited to, atrial fibrillation, diabetes, long QT syndrome, and hypertrophic cardiomyopathy.
The compound 4-(pyrimidin-2-ylmethyl)-7-(4-(trifluoromethoxy)phenyl)-3,4-dihydrobenzo[f][1,4]oxazepin-5(2H)-one, designated herein as Compound I, is known to be a selective late sodium current inhibitor, as described for example in WO 2013/006485.

Abstract: Also provided are processes of manufacture and methods of using the crystalline forms.


An approach for authenticating parties engaged in a web-based transaction without compromising the integrity or anonymity of the parties is described. An authentication platform receives, from a first application associated with a first party, an authentication request that has been redirected in response to a transaction initiated with a second application associated with a second party. The authentication platform forwards knowledge based assessment information to the first application that is based on determined behavioral information for authenticating the second party to the first party. A valid response to the knowledge based assessment information by the first application provides authentication of the first party to the second party.

It is commonplace for one party to engage in a transaction with another party directly over a public data communications network, such as the global Internet. In these types of transactions, be it for engaging in a purchase transaction, facilitating a trade or exchanging vital information, the parties involved may be anonymous to one another. The anonymity associated with parties corresponding via the Internet or other network medium can lead to trust concerns that prevent otherwise successful transactions from occurring, particularly in the area of commercial collaboration. Unfortunately, many traditional methods for authenticating parties involved in a transaction require the exchange of sensitive information regarding the parties involved, which jeopardizes anonymity.
Based on the foregoing, there is a need for an effective and reliable way to authenticate parties without compromising the integrity or anonymity of respective members.

Method and system for providing behavioral bi-directional authentication


A drill guide, for use in the preparation of a dental implant site, having a tubular body with a screw thread along at least a section thereof. The drill guide has a stent with screw threads that are complementary screw threads allow the removable attachment of the tubular body to the stent and allows the axial advancement of the tubular body toward the bone. The complementary screw threads are provided on a matrix material of the stent.

Dental implants are used to secure a false tooth or for securing other dental items, such as bridges, in the oral cavity. The aim is to achieve a firm, durable, yet flexible, intra-oral connection to the jaw bone of screw-type or cylinder-type structures made of a biologically compatible material. Such implants are often made of titanium, an alloy thereof, and some have been sprayed with titanium or coated with hydroxyl apatite.
Various forms of dental implants are described in Israel Patents 92360, 92561 and 92564. These specifications do not relate to the problem of drilling an accurate preparatory hole for the insertion therein of the implant.
Briefly, dental implants are inserted in the jaw bone of a patient using the procedure described below:
1. Clinical and radiographic examination. PA0 2. Preparation of a stone model of the relevant existing structures. PA0 3. Preparation of a transparent acrylic stent for use in the mouth during surgery. PA0 4. Bone exposure by incision and dissection of soft tissues. PA0 5. Drilling the bone, using clear acrylic splint as an (imprecise) guide. PA0 6. Hole expansion by further drilling or dilation. PA0 7. Insertion of a screw or cylinder implant. PA0 8. Adding a cover screw or healing cap to the implant head. PA0 9. Stabilizing the tissues to either cover the implant, or to leave the implant cover flush with the gingival tissue. PA0 10. Osseointegration, which takes place naturally during the following 4-12 months. PA0 11. Implant re-exposure, if necessary, and subsequent construction of a dental prosthesis on the implant infrastructure.
When the bone is of small dimensions or is porous, special techniques are used which include bore expansion by hammering in a dilator, increasing bone dimensions by making a chisel fracture, and floor elevation of the maxillary sinus cavity (bone grafting) for increasing the available depth and/or increasing the width by onlay bone or bone-substitute grafting.
Drilling of the bone is carried out by one of the following known techniques: a) free-hand drilling; b) drilling through a prepared hole in a plastic matrix, or c) drilling through a smooth, non-threaded, cylindrical metal drill bush held in a plastic matrix.
The initial hole in the jaw bone must be drilled at the correct location, typically half-way between two existing teeth, at the correct angle in all directions, and to the required depth. A faulty drilling trajectory is very difficult to correct and can cause damage to nearby structures such as the inferior alveolar nerve; it can also cause pain to the patient and, in some cases, cause implant failures.
Free-hand drilling is a technique which is suitable only for the exceptionally skilled and experienced dental surgeon. There will be difficulties when the bone being drilled is not homogeneous, the drill will be deflected and a too-large diameter hole will result.
Technique b) and, to an even greater extent, technique c), greatly reduce the above-mentioned risks associated with free-hand drilling. However, where the guide bore is directly drilled in the plastic matrix as in technique b), play between the bore and the drill is likely to develop, resulting in imperfect results. This problem is overcome in technique c).
The use of a metal drill bush has long been known in production engineering, and standard sizes are commercially available `off the shelf`. [See Wilson, Handbook of Fixture Design, McGraw-Hill Publishers, pp. 15-20 (1962)]. In technique c), this time-tested method is adapted to dental requirements.
However, the non-threaded metal drill bush is difficult to remove from the stent if it needs to be removed and redirected in a corrected drilling axis; furthermore, it is not axially adjustable. Lack of axial adjustment is a disadvantage when drill stability needs improvement, such as when the inner extremity of the bush is spaced from the bone being drilled. In addition, there is no secure means for adding attachments to the bush, such attachments having many important uses, as will be explained further below.
It is therefore an object of the present invention to obviate the disadvantages of the prior art dental implant drilling devices, and to provide a drill guide which is readily removable and replaceable, and which can be prechecked for accuracy.
It is a further object of the present invention to provide a drill guide which is readily depth-adjustable.
It is a still further object of the present invention to provide means for the secure attachment to the drill guide of attachments such as a bush for drilling a pilot hole, holders and guides for dilators and chisel heads, a radiographic marker, a depth gauge, sinus lift attachments, and the like.
Further objects of the invention will become apparent from the following description.