Naphtha is selectively hydrodesulfurized using deactivated hydrotreating catalyst to remove sulfur while minimizing loss in octane level due to olefin saturation.

1. Field of the Invention
This invention concerns a process for removing sulfur from naphtha, a petroleum product used to make fuels like gasoline. Specifically, deactivated hydrotreating catalyst is used to selectively hydrodesulfurize naphtha while minimizing olefin loss.
2. Description of Related Information
It is well known that air pollution is a serious environmental problem. A major source of air pollution worldwide is the exhaust from hundreds of millions of motor vehicles due to fuel combustion. Laws and regulations have been enacted reflecting the need to reduce harmful motor vehicle emissions through more restrictive fuel standards. Fuels containing sulfur produce sulfur dioxide and other pollutants leading to a host of environmental concerns, such as smog and related health issues, acid rain leading to deforestation, water pollution, as well as other environmental problems. To help reduce or eliminate these environmental problems, the sulfur content of fuels has been, and will continue to be, restricted to increasingly smaller concentrations, such as less than 100 or even 50 parts per million (ppm).
The problem of sulfur in fuels is compounded in many areas where there is diminishing or no domestic source of crude oil having relatively low sulfur content. For example, in the United States the supply of domestic oil production relies increasingly on lower grade crude oil with higher sulfur content. The need for lower sulfur content fuel therefore increases demand for imported oil having lower sulfur content increasing trade imbalance and vulnerability due to dependence on foreign sources of oil.
The sulfur content in crude oil can take the form of a wide variety of both aliphatic and aromatic sulfurous hydrocarbons. Various techniques have been developed for removing sulfur compounds. One such technique, called hydrodesulfurization (HDS), involves catalytically reacting hydrogen with the sulfur compounds. The general HDS reaction is illustrated in Equation 1. EQU RSR'+H.sub.2 .fwdarw.RH+R'H+H.sub.2 S EQU Equation 1: Hydrodesulfurization Reaction
In Equation 1, the sulfur compound, RSR', may be: a thiol or mercaptan, where R is hydrocarbyl and R' is hydrogen; a sulfide or disulfide, where the sulfur is connected to another sulfur atom in R or R' hydrocarbyl groups; or may be a thiophene where R and R' are connected to form a heterocyclic ring. The HDS reaction consumes hydrogen (H.sub.2) and produces hydrogen sulfide (H.sub.2 S) and hydrocarbons wherein the sulfur atom is replaced by two hydrogen atoms. The hydrogen sulfide can then be separated to give a petroleum product in which the sulfur is significantly reduced or substantially eliminated.
HDS is one process within a class of processes called hydrotreating, or hydroprocessing, involving the introduction and reaction of hydrogen with various hydrocarbonaceous compounds. General hydrotreating reactions with oxygen compounds, nitrogen compounds and unsaturated hydrocarbons, including olefins, are illustrated in Equations 2, 3 and 4, respectively. EQU ROR'+H.sub.2 .fwdarw.RH+R'H+H.sub.2 O EQU Equation 2: Hydrodeoxygenation Reaction ##STR1##
The hydrotreating reactions can occur simultaneously to various degrees when sulfur-, oxygen-, nitrogen-containing and unsaturated compounds are present in the petroleum. The hydrotreating reactions are exothermic, producing heat. Such hydrotreatment has been used to remove not only sulfur, but to also remove nitrogen and other materials, like metals, not only for environmental considerations but for other uses, such as to protect catalysts used in subsequent processing from being poisoned by such elements. See, for example, Applied Industrial Catalysis, Volume I, edited by B. E. Leach, Academic Press (1983); Chemistry of Catalytic Processes, by B. C. Gates et al., McGraw-Hill (1979); and Applied Heterogeneous Catalysis: Design Manufacture Use of Solid Catalysts, by J. F. LePage et al., Technip, Paris (1987).
Olefins are useful in fuel feedstock by raising the octane number of the fuel, increasing its value and performance properties. For example, cracked naphtha typically contains over 20 weight percent olefins having octane numbers that are higher than the corresponding saturated hydrocarbons. HDS of naphtha using standard hydrotreating catalysts under conditions required for sulfur removal produces a significant loss of olefins through hydrogenation. This produces a lower grade fuel which then needs more refining, such as isomerization, blending, or other refining, to produce higher octane fuel, adding significantly to production expenses.
Selective HDS to remove sulfur while minimizing hydrogenation of olefins and octane reduction by various techniques, such as selective catalysts, have been described. For example, U.S. Pat. Nos. 4,132,632 (Yu et al.) and No. 4,140,626 (Bertolacini et al.) disclose selective desulfurization of cracked naphthas by using specific catalyst having particular amounts of Group VIB and VIII metals on magnesia support. U.S. Pat. No. 4,149,965 (Pine et al.) discloses a process for starting-up naphtha HDS using partially deactivated hydrotreating catalyst under relatively low pressure of less than 200 psig. The catalyst is partially deactivated using a substantially non-metals containing, hydrocarbonaceous oil for a time ranging from about 10 hours to about 20 days. U.S. Pat. No. 2,983,669 (Noll) discloses processes for treating petroleum with high sulfur content using fractionation and HDS. Noll suggests that the HDS catalyst may be one which does not readily promote hydrogenation, such as a partially spent catalyst.
Hydrotreating catalysts age, losing activity during use by collecting deposits of carbonaceous material and/or impurities, such as metals, from the treated feedstock. Eventually, with increased deposition the catalyst is no longer able to provide effective hydrotreating. The deactivated catalyst may be regenerated. The regenerated catalyst can be reused but is generally less effective than fresh catalyst by requiring higher temperature to give the desired activity and becoming deactivated more quickly than fresh catalyst. Although hydrotreating catalysts can usually be repetitively regenerated, they eventually become irreversibly deactivated, or spent, essentially losing their intended hydrotreating utility.
Spent hydrotreating catalysts have been used in hydrotreating, including HDS. For example, U.S. Pat. No. 3,876,532 (Plundo et al.) discloses a process for hydrotreating middle distillate, virgin oils using spent hydrotreating catalyst under extremely mild conditions to reduce acid and mercaptan content, to remove sulfur below 0.2 weight percent, or 2,000 ppm. U.S. Pat. No. 4,414,102 (Rankel et al.) discloses the use of spent HDS catalyst to transform nitrogen- or oxygen-containing compounds to sulfur-containing compounds followed by mild HDS treatment.
It would be desirable to have a process for removing sulfur from fuel feedstocks, like naphtha, containing olefins which minimizes loss of octane value using an inexpensive procedure under a wide range of conditions, to provide a cleaner environment along with a more stable economy.

Selective hydrodesulfurization of naphtha using deactivated hydrotreating catalyst


A method and apparatus for bonding integrated circuits uniquely suited to high volume tag production is described, where conductive material of a substrate at the die-attach-area is cut before an IC chip or transponder is placed on the conductive material over the cut and bonded. The apparatus performs the method of placing a first chip on a substrate having a conductive layer, measuring the location of the first chip on the substrate, cutting the conductive layer at a location of an expected subsequently placed chip to form a cut based on the measured location of the first chip, and placing the subsequently placed chip on the substrate over the cut.

Chip bonding is costly. The two largest components of the cost of RFID tags today are the integrated circuit and the attachment of that circuit to an antenna structure. Moore's law and increasing volume are helping to drive the IC cost down, but bonding is a mechanical process and does not benefit from the same technology advances or economies of scale.
Current methods of chip bonding do not adequately address cost. A two-step approach of an intermediary "strap" achieves incremental cost improvement by relocating the costs. However, straps do not address the problem directly, as bonding is still required, but to a smaller tag. Moreover, straps add another step to bond the strap to the big tag.
Current manufacturers using standard bonding technology with straps, want straps to be like traditional bonding surfaces, that is, hard and inflexible. But such straps do not lend themselves to easy integration into squishy flexible tags. The known standard bonding processes are all strap-based solutions, and therefore less than ideal.
One related art bonding method, called fluid self-assembly, provides insufficiently robust bonds. Because the chips find their own way into bonding sockets, the chips cannot use any adhesives or flux, since anything sticky prevents free motion of the chip into the sockets. Then the bond is made at a tangent between the chip bonding pad and the side of the bonding cavity. This flat-to-edge bond is different than and less reliable than traditional bonds, which are made flat-to-flat. As an analogous example illustrating problems inherent to flat-to-edge bonds, consider trying to stand a playing card on edge, rather than laying it flat on a table. Fluid self-assembly also places restrictions on the type of substrate that can be used. This may not be a problem for just making straps; but it is certainly a problem for foregoing the strap and putting the chip right on the tag.
A known bonding process is a gallant brute force attempt to make standard bonding faster. Instead of having one vacuum head pick up one chip and place it on one strap, a plurality of heads (e.g., 60) in lock step picks up that number (e.g., 60) of chips and puts them on that number (e.g., 60) of straps. This process suffers from the problem of keeping all of the number (e.g., 60) of chips aligned correctly at the same time.
Bonding RFID chips is more like processing diodes and resistors than processing other kinds of chips. One new RFID strap line uses a traditional tape automated bonding process, with a sprocket-fed 35 mm tape of hard straps inching its way through a traditional flip-chip placement and bonding head. At 4.75 mm pitch, four lanes wide, and 10,000 chip bonds per hour, their tape advances through the bonding procedure at about 0.65 feet per minute. It would be beneficial if a chip bonding process could produce more bonded chips in less time.
In order to consider why the art has not bonded chips as exemplified by the preferred embodiments of the invention discussed below, it may help to compare standard electronic chip components to RFID tags. Standard electronic chip components are known and generally found on printed circuit boards. A bare IC is bonded to a carrier by wire bonding or flip chip. Then a package is molded around the carrier and chip. The package is then put onto a printed circuit board via thru-hole or surface mount assembly. In summary, typical standard chip components: need to be compatible with multiple PCB assembly technologies, including solder baths, solder waves, IR reflow, and a variety of cleaning and baking steps; want more and more computational power put in single chip assemblies; and are made to last. In contradistinction, RFID tags: are never soldered or baked or cleaned; are complete unto themselves and do not have to be integrated into any other system; want the bare minimum computational power to minimize cost and energy consumption (which translates into read distance); and do not face the same power dissipation or environmental requirements as standard chips.
To meet their design requirements, standard chip assemblies usually start with relatively stiff and heavy substrates, at least compared to RFID tags. Ceramics and fiberglass are common. These are meant to be tough and resistant to thermal influences. Usually the standard chip substrates are etched. Laser cutting is expensive because the standard chip substrates are thick and have high thermal masses.
RFID tags are substantially different. The metal layer is thin and flexible (or non-rigid) by comparison. The back or substrate of each tag is soft polypropylene or paper. The substrates are easily to punch, cut, dimple, and weld. The preferred embodiments of the invention reinvent bonding taking advantage of these different properties.
A known wire bonding process is disclosed in U.S. Pat. No. 5,708,419 to Isaacson, et al., the contents of which are incorporated by reference herein in its entirety. Isaacson discusses the bonding of an IC to a flexible or non-rigid substrate which generally can not be subjected to high temperatures, such as the temperature required for performing soldering processes. In this wire bonding process, a chip or die is attached to a substrate or carrier with conductive wires. The chip is attached to the substrate with the chip front-side face up. Conductive wires are bonded first to the chip, then looped and bound to the substrate. The steps of a typical wire bonding process include:                1. advancing web to the next bond site        2. stopping        3. taking a digital photograph of the bond site        4. computing bond location        5. picking up a chip        6. moving the chip to the bond site        7. using photo feedback to adjust placement to the actual site location        8. placing or depositing chip        9. photographing the chip to locate the bond pads        10. moving the head to the chip bond pad        11. pressing down, vibrating and welding conductive wire to the bond pad        12. pulling up and moving the chip to the substrate bond pad, trailing wire back to the chip bond        13. pressing down and welding that bond        14. pulling up and cutting off the wire; and        15. repeating steps 10-14 for each connection
In contrast, the interconnection between the chip and substrate in flip-chip packaging is made through conductive bumps of solder that are placed directly on the chip surface. The bumped chip is then flipped over and placed face down, with the bumps electrically connecting to the substrate.
Flip chip bonding, a current state of the art process, is expensive because of the need to match each chip to a tiny, precision-cut bonding site. As chips get smaller, it becomes even harder to precisely cut the bonding site. However, the flip-chip bonding process is a considerable advancement over wire bonding. The steps of a typical flip-chip bonding process include:                1. advancing web to the next bond site        2. stopping        3. photographing the bond site        4. computing the bond location        5. picking up the chip        6. moving the chip to the bond site        7. using photo feedback to adjust placement at the actual site location        8. placing the chip        9. ultrasonically vibrating the placement head to weld chip in place; and        10. retracting the placement head
Steps 1 through 8 of each of the above bonding processes are substantially the same. The web must stop to locate the conductive gap in the substrate and precisely place the IC. The related art processes require that the web is stopped and measured (e.g., photographing the bond site, containing the bond location, using photo feedback to adjust placement at the actual site location) so that the chip can be accurately placed as desired adjacent the gap and bonded.
In designing an efficient chip placement process that can be integrated into RFID tags, the inventors discovered that it is beneficial to avoid anything that is not consistent with a continuous rolling printing press. Stopping and starting the line always slows things down. It would be beneficial to adjust tooling to operate on a chip that is continuously advancing down the line at a known rate of travel.
Retracing a path during the bonding process takes time, causes vibration, and wears mechanical linkages. These linkages also create uncertainty in absolute position. Rotating or continuous devices are thus preferred over reciprocating devices.
The greater the number of mechanical connections in a bonding process, the less certainty there is in precise position. Every jointed or flexible linkage introduces a certain amount of randomness as the web and chips wiggle around. IC dimensions are tiny. It does not take a lot of mechanical links to move chip placement out of critical alignment.
With security tags, you cannot rely on any precise dimension set previously. The relative position of things varies across the web, from one end of the roll to another, from place to place, and from time to time. That is simply the reality of working with inexpensive materials. For IC bonding processes, the manufacturer must constantly adapt to how the material is really behaving, rather than counting on it to behave as intended.

Die attach area cut-on-fly method and apparatus


The present invention provides a method for directly producing lactide by subjecting lactic acid to a dehydration reaction in the presence of a catalyst comprising a tin compound, preferably, a tin (IV) compound, wherein lactide can be produced directly or by one step from lactic acid, without going through the step of producing or separating lactic acid oligomer. The method of the present invention has advantages of causing no loss of lactic acid, having a high conversion ratio to lactic acid and a high selectivity to optically pure lactide, and maintaining a long life time of the catalyst. Further, since lactic acid oligomer is not or hardly generated and the selectivity of meso-lactide is low, the method also has an advantage that the cost for removing or purifying this can be saved.

Polylactide (PLA) is a polymeric material having optical characteristics, which has been used as a biodegradable or bioabsorbable medical material for surgical suture thread, microcapsules for injection or the like, and recently, is employed as an environmental-friendly biodegradable material which can be used in the production of various high-molecular products such as packaging materials, home appliances, office supplies, vehicle interior materials, or the like.
In order for polylactide to be used for the above usage, it is required to have a high optical purity (D-type or L-type optical isomer) and have high molecular weight. For this, lactide, which is a monomer used for producing polylactide, should have high optical purity and chemical purity.
The following Reaction Chart 1 shows routes for producing lactide from lactic acid or lactic acid ester. In case of using lactic acid as a starting material, lactide can be produced through a course of Route-1/Depolymerization or a course of Route-4. In case of using lactic acid ester as a starting material, lactic can be produced through a course of Route-2 or a course of Route-3/Depolymerization.

The course of Route-1/Depolymerization is a conventional method for producing lactide by using lactic acid as a starting material, wherein lactide is obtained via a two-step process consisting of firstly polymerizing lactic acid under reduced pressure to produce an oligomer or prepolymer in a molecular weight of about 500˜5,000, and then depolymerizing the produced prepolymer with the flow of inert gas under a reduced pressure (U.S. Pat. No. 5,274,073; U.S. Pat. No. 5,247,059; U.S. Pat. No. 5,274,127; U.S. Pat. No. 6,277,951). NatureWorks (a company in USA) commercially produces L-type PLA bioplastic in a scale of about 140 thousand tons a year by a manner consisting of producing lactide from L-type lactic acid via said two-step process, and then subjecting the resulting lactide to a condensation polymerization reaction in a polymerization reactor. Recently, Purac (a company in Netherlands) has also built a plant in Thailand and started a commercial production of lactide via the above two-step reactions.
However, in case of going through the two-step process, the products such as prepolymer and lactide are deteriorated when they reside for a long period of time within the reactor, and thus there is a problem that the amount of byproducts produced such as meso-lactide increases during the production of L-type or D-type lactide. The two-step process of Prepolymerization-Depolymerization of lactic acid has problems that part of the prepolymer or oligomer which is the reactant in the Depolymerization step is not decomposed to lactide and further polymerizes or forms a waste material mixed with the catalyst component included in the prepolymer. Also, the two-step process involves problems that a vacuum pump for making a high vacuum is required and the reaction device is complex to make the device cost high.
The course of Route-2 or Route-3/Depolymerization is a method using lactic acid ester as a starting material. There are several references disclosing a process via Route-2 for the production of lactide directly or by one-step from lactic acid ester, or a process via Route-3 and Depolymerization for the production of lactide through prepolymer from lactic acid (e.g. Korean Patent Application No. 10-2009-0043985, Japanese Patent Laid-Open No. 1999-036366, Japanese Patent Laid-Open No. 1993-286966, and Japanese Patent Laid-Open No. 1994-031175).
Finally, the process via Route-4 relates to a method for producing lactide directly or by one step from lactic acid as a starting material. A method for the production of lactide by this process or a catalyst for the production thereof has ever been suggested, but the suggested method and catalyst have failed to provide satisfactory results in conversion ratio and selectivity when considering the process efficiency.
U.S. Pat. No. 5,332,839 suggested a direct production of lactide by a process via Route-4. For example, it reports that lactide can be produced directly from lactic acid without going through a prepolymer by gasifying lactic acid at a high temperature of at least 200° C. and then reacting it in a fixed layer charged with a solid acid such as Al2O3. In case of using a solid acid catalyst such as Al2O3, however, there are limitations such that the productivity of lactic acid oligomer is high, the yield of lactide is low, carbon monoxide is generated by the degradation of lactic acid during the reaction, and the life of the catalyst is short.
U.S. Pat. No. 5,138,074 also suggested a direct production of lactide by a process via Route-4. For example, it discloses an example for producing lactide directly from lactic acid in the presence of SnO. However, it does not disclose results such as conversion ratio of lactic acid, selectivity of lactide, production ratio of oligomer byproduct, and thus it is difficult to refer to the catalyst and reaction in detail.
Many methods for producing lactide by using lactic acid ester as a starting material (Route-2 and Route-3) and methods for producing lactide by a two-step process by using lactic acid as a starting material (Route-1 and Depolymerization) have been developed. The method for producing lactide directly or by one step from lactic acid as a starting material (Route-4) has not been employed in a commercial process due to the absence of a suitable catalyst, although it has an advantage that the process is simple when compared with the two-step process by the Prepolymerization/Depolymerization of lactic acid (Route-1 and Depolymerization).

Method for producing lactide directly from lactic acid and a catalyst used therein

KOREA RESEARCH INSTITUTE OF CHEMICAL TECHNOLOGY


A two-step method is provided for the production of a pressure sensitive adhesive by radiation curing. The method comprises the steps of initially irradiating a monomer composition comprised of 65 to 100 percent by weight of monomers having a Tg<0.degree. C. to form a coatable syrup, followed by adding at least one monomer having a Tg>0.degree. C. and at least one multifunctional monomer or oligomer to said syrup, and further irradiating the mixture to form a pressure sensitive adhesive.

The present invention is directed to a two-step method for the production of a pressure sensitive adhesive by radiation curing.
It is well known that pressure sensitive adhesives may be prepared by radiation curing of monomer admixtures. See, for example, U.S. Pat. Nos. 4,181,752; 4,379,201; 4,421,822; 4,513,039; 4,522,870; 4,587,313; 4,665,106; 5,183,833; 4,737,559; 5,302,629; 5,462,977; 5,536,759; 5,552,451; 5,618,899 and 5,683,798. Indeed, attempts have been made to provide reinforced pressure sensitive adhesives by incorporating reinforcing macromers in the radiation curable composition. See, for example, PCT/US97/09601 and WO94/13750.
It is, however, desirable to provide a method for the production of a macromer-free pressure sensitive adhesive by radiation curing which does not require the presence of a reinforcing macromer to yield a pressure sensitive adhesive exhibiting desirable adhesive properties.

Two-step method for the production of pressure sensitive adhesive by radiation curing


A system for detecting precursor seismic electromagnetic waveforms includes a transducer connected to the emitter of a transistor and arranged to conduct in response to an electromagnetic seismic precursor. Conduction by the transducer initiates a current in the transistor's base-emitter junction, reversing the bias of the function and enabling a resonating current in the junction.

The present invention relates to an electromagnetic detector system, and in particular a seismic electromagnetic detector system suitable to detect precursor electromagnetic waveforms of earthquakes.
Often a devastating toll in life and property is taken by earthquakes. It has long been recognized that the toll could be reduced if people within an impending quake's focal area were warned to prepare. Although preparations are unlikely to prevent structural damage to commercial and residential buildings, or to infrastructure such as bridges and roadways, preparation could reduce deaths and serious injuries by people seeking appropriate shelter or retreating from dangerous locations, such as unreinforced brick buildings. Moreover, preparation is likely to reduce the psychological trauma often attributable to an earthquake's sudden onslaught. In addition, preparation will likely reduce personal property causality, such as that related to structural or utility failures, and the fires often associated therewith. Accordingly, it is desirable to forecast the occurrence of earthquakes.
An earthquake's toll results from the seismic waves defining the earth-quake. Seismic waves include two types: body waves and surface waves. Body waves comprise primary (or P) waves and secondary (or S) waves that propagate within the earth's body. P waves are longitudinal waves that alternately push (compress) and pull (dilate) the ground in the direction of propagation. S waves are transverse waves that shear the ground in planes perpendicular to the direction of propagation.
Surface waves comprise Love waves and Rayleigh waves that propagate at or near the earth's surface. Love waves shear the ground sideways at right angles to the direction of propagation, much like S waves, but without S waves' vertical shearing. Rayleigh waves displace the earth both vertically and horizontally will travel an elliptical path as the wave passes, much like a water molecule in rolling ocean waves.
Body waves travel more rapidly than surface waves. Of the surface waves, Love waves generally travel faster than Rayleigh waves and, of the body waves, P waves generally travel faster than S waves. When an earthquake is occurring, the P waves are felt first, like a thud or blow, and thereafter the S waves arrive, as indicated by up-and-down and side-to-side motion. Thereafter, the surface waves strike, causing the ground to shake side-to-side and to roll.
The body and surface waves generally are monitored during an earthquake to gauge the earthquake's intensity. Being contemporaneous with and defining earthquakes, these waves cannot be used for forecasting. Forecasting relies on identification of other physical parameters that, in their occurrence or variance, indicate an impending earthquake. These parameters, when indicating an impending earthquake, are sometimes referred to herein as "precursor seismic activity."
As reported by an article written by Evelyn Roeloffs, a team in 1989 from Stanford University happened to be listening for low-frequency magnetic noise in the Santa Cruz Mountains south of San Francisco with a large coil. On September 12, they noticed an unusual signal with a period between 5 and 20 seconds, which was followed by a background noise increase October 5. On October 17, the background noise rose to a high level, and three hours later the magnitude 7.1 Loma Prieta earthquake hit, rocking the San Francisco Bay area. The earthquake was centered less than five miles from their measurement coil. Since then, one explanation after another for these noise increases have been eliminated, leaving open the possibility that they truly were earthquake precursors. Hoping to repeat the experiment, they have deployed three similar instruments in California, two of them near Parkfield, which is a site of intensive earthquake prediction research. Electromagnetic precursors are a controversial subject worldwide. An international group of scientists continues to deliberate whether "seismic electric signals" recorded in Greece are precursors to earthquakes. But there is growing consensus that the Earth's electrical resistance decreased in association with the 1976 Tangshan, China, earthquake. The article goes on to state that as with all potential earthquake prediction techniques, there won't be significant progress until a good-sized earthquake happens in a closely monitored location.
It would appear that the Stanford detector may have detected precursors to the Loma Prieta earthquake so its design has been the basis for further earthquake prediction devices. The Stanford detector is a large coil of wire designed to be resonant in the range of approximately 1-30 hertz. To obtain sensitivity within this range the number of coils and the core for the coils are both selected accordingly. Unfortunately, such a coil is also sensitive to man-made noise, lightning, and electromagnetic fields from storms and atmosphere. Accordingly, the Stanford detector needs to be located in remote areas to minimize the detection of extraneous noise. Furthermore, other researchers have had difficulties using the Stanford detector to detect any precursor activity to earthquakes, let alone actually predict earthquakes, so it would appear that the Stanford detector would need to be located near the epicenter of a large earthquake to be effective.
Varotsos et al., U.S. Pat. No. 4,612,506, disclose a method of forecasting earthquakes as a function of transient variations in electric earth current. Varotsos discovered that the electric currents which normally flow in the earth, termed telluric currents, undergo transient changes or variations of a specific nature or character at times before the occurrence of an earthquake. Specifically, Varotsos found that earthquakes are preceded by a first transient variation in telluric earth current measurable as a voltage on the order of hundreds of microvolts per earth-meter having about one minute duration that occurs from six to eight hours before the quake, and a second transient change in earth current measurable as a voltage on the order of tens of millivolts per earth-meter having a duration of a few milliseconds and occurring between thirty seconds and four minutes before the quake. While such a method is useful for prediction of impending earthquakes, the detection of precursor electromagnetic seismic activity only up to eight hours before a quake is an inadequate length of time to warn the public.
Varotsos' detection system involves the simultaneous measurement of pre-earthquake long waves or earth currents at a number of points in the earth by using multiple elongated conductive cylinders. More specifically, the pre-earthquake long waves or earth currents are measured simultaneously at two or more points on the earth surface. Each of the cylinders measures a transitory current that propagates in the crust. The distance that the waves travel in the earth between cylinders results in a small voltage drop of the waves between them. An operational amplifier produces a signal which is the differential potential between the cylinders. The point of origin of the electromagnetic waves or earth currents is computed from the amplitude ratio of the detected signals. As described in Varotsos the cylinders are vertically aligned with axes in orthogonal planes to the direction to the epicenter. Presumably, the theory is that the electromagnetic waves from seismic activity propagate radially outwardly from the epicenter of the pending quake striking the cylinders. With axes in orthogonal planes, the cylinders are oriented to expose the maximum surface area in a direction normal to the epicenter in order to maximize the detected potential difference. Unfortunately, Varotsos' system has noted that periodically earthquakes occur where there were no electromagnetic precursors detected prior to the earthquake.
For the system taught by Varotsos, the changes in earth current preceding an earthquake of a given intensity must be determined empirically for each location because the intensity changes in the earth current are a result of the distance from the earthquake's epicenter, earth conductivity, and the magnitude of the quake itself, all of which vary from location to location. This uncertainty is unacceptable for use in areas that do not experience frequent earthquakes suitable to calibrate the detector. For regions that have infrequent but devastating earthquakes, it would take several disasters to calibrate the detector. In addition, it is difficult to predict earthquakes if the detectors are located on opposing sides of a subterranean feature, such as a fault line.
Tate et al., U.S. Pat. No. 4,628,299 disclose, a seismic warning system using a radio frequency energy monitor. However, such a system is not accurate because changes in the tides and other factors influence the radio frequency field strength which requires statistical calculations for which it attempts to compensate. Accordingly, it is not feasible to discriminate small seismic activity.
Takahashi, U.S. Pat. No. 5,904,943, discloses a system similar to Varotsos et al. that includes a three dimensional distribution of the sources and intensities of the long waves or earth currents, and predicts the focal region, scale and time of occurrence of earthquakes. Takahashi defines the long waves and earth currents as sinusoidal measurements having a frequency not exceeding 300 kHz. With the premise that the waves are sinusoidal in nature, Takahashi teaches the use of an antenna for the long waves and measuring the voltage difference between a pair of vertical cylinders for the earth currents. Further, Takahashi teaches that electromagnetic waves over 300 kHz are attenuated at 0.01-1.0 dB/m and therefore seldom observed near the earth surface so no attempt is made to detect waves with higher frequencies. Accordingly, the detectors (cylinders or antennas) are located between 100 km and 500 km from the epicenter depending on the earthquake size to be detected. Takahashi's method includes many of the problems associated with the method taught by Varotsos.
Takahashi proposes that an earthquake is a sudden shifting of the earth's crust along a fault plane that occurs when the stress within the crust comes to exceed the deformation limit. However, as the rock forming the crust of the earth is not homogeneous, small-scale disintegration and dislocation of the rock occurs locally at points within the fault plane destined to become the focal region before the earthquake actually occurs. This gives rise to electromagnetic waves (long waves and earth currents). As a result, Takahashi teaches that it becomes possible to predict the occurrence of an earthquake from the electromagnetic waves from the wave source region. However, this method is further limited to about one week prior to the earthquake.
Helms, U.S. Pat. No. 4,507,611, discloses a method of detecting surface and subsurface anomalies of the earth using vertical current measurements. The vertical current manifests itself as alternating current signals which can be measured and represents surface and subsurface anomalies. Local variation in the detected current signals are measured and correlated with the spatial relation to the points of measurement to determine significant measurements indicative of surface and subterranean anomalies.
Weischedel, U.S. Pat. No. 4,219,804, teaches a circuit suitable for identifying electromagnetic radiation signals caused by nuclear detonations and discriminating against false indications by lightning. The circuit is designed to detect the electromagnetic waveform from the nuclear detonation or lightning strike, as a single negative-going waveform. In order to detect the signal an antenna is used. If a signal is detected that exceeds a threshold magnitude, then three circuits are enabled: a zero crossover discriminator, a rise time discriminator, and a precursor discriminator. The zero threshold detector determines whether the signal crosses the zero axis within a required time period, and if this occurs, it provides an output pulse to an "AND" gate. The rise time discriminator determines whether the signal rises to its peak within a required time period, and if this occurs, it provides an output pulse to the "AND" gate. The precursor discriminator compares the peak amplitude detected with a previous signal detected, if any, occurring within a certain previous time. If the previous signal was detected within the prescribed time period then no signal will be provided to the "AND" gate, indicative that this was lightning. If no previous signal was detected within the prescribed time period then a signal is provided to the "AND" gate, indicative of the possibility of a nuclear detonation. The circuit of Weischedel is designed to detect nuclear detonations by sensing a single waveform and discriminate this against a prior waveform within a prescribed time period in order to discriminate against lightning. While suitable for lightning and nuclear detonations, this circuit is not capable of predicting earthquakes.
There are numerous systems for geophysical exploration of the earth primarily for the detection of mineral and oil deposits. These systems are principally based on the detection and analysis of alternating currents, such as generally sinusoidal currents, from within the earth. Some systems impose a waveform into the earth, while others detect changes in existing earth currents. Such systems include, for example; Nilsson, U.S. Pat. No. 3,701,940; Miller, et al. U.S. Pat. No. 4,041,372; T. R. Madden et al., U.S. Pat. No. 3,525,037; Hearn, U.S. Pat. No. 3,976,937; Barringer, U.S. Pat. No. 3,763,419; L. B. Slichter, U.S. Pat. No. 3,136,943; G. H. McLaughlin et al., U.S. Pat. No. 3,126,510; and Weber, U.S. Pat. No. 4,044,299.


A method for manufacturing an electrode for an all solid battery including the steps of coating a current collector with a slurry including an active material, a conductive material, and a polyimide-based binder; and melting a solid electrolyte having a melting temperature of 50° C. to 500° C. and applying it onto the coating layer and an electrode manufactured therefrom.

From the viewpoints of capacity, safety, power output, enlargement, miniaturization and the like of the battery, various batteries that can overcome the limitations of lithium secondary batteries are currently being studied.
Representatively, a metal-air battery, which has very large theoretical capacity compared to the lithium secondary battery, in terms of capacity, an all solid battery with no risk of explosion in terms of safety, a supercapacitor in terms of output, a NaS battery or RFB (redox flow battery) in terms of enlargement, a thin film battery in terms of miniaturization, etc., are continuous studied in academia and industry.
Among these, the all solid battery is a battery in which a liquid electrolyte used in a conventional lithium secondary battery is replaced with a solid electrolyte, and the all solid battery can significantly improve safety since a flammable solvent is not used in the battery and thus the ignition or explosion due to the decomposition reaction of the conventional electrolyte solution is not occurred. In addition, there is an advantage that since Li metal or Li alloy can be used as the material of the negative electrode, the energy density with respect to the mass and volume of the battery can be remarkably improved.
There are roughly two methods that can be used in manufacturing the all solid battery.
First, there is a method of pressurizing after dry-mixing. Specifically, after adding the solid electrolyte powder to the active material powder and conductive material powder and then dry-mixing, the electrode composite powder is pressurized onto the current collector to form the electrode.
In addition, there is a method to pressurize after wet-mixing. Specifically, after wet-mixing the active material powder and the conductive material powder with the solid electrolyte powder and the binder solution, the electrode material is coated on the current collector in the form of a slurry, dried and pressurized to form the electrode.
Among them, in the case of the dry-mixing method, there is a disadvantage that the interfacial resistance between the current collector/electrode/electrolyte is high and the control of the pores in the electrode is difficult and there is also a disadvantage that the contact between the electrode material and the electrolyte is not maintained because there is no binder.
In addition, in the case of the wet-mixing method, there are disadvantages that it is difficult to apply the binder and the solvent used in the conventional secondary battery process, and when the solvent is vaporized, the pores are generated in the electrode and thus the resistance in the electrode/electrolyte is increased.
Additionally, in both of the above methods, the three kinds of electrode materials of active material, electrolyte, and conductive material must be evenly dispersed, but there is a high possibility that the particles are not well dispersed and thus there is a disadvantage that the resistance of the electrode is large. Further, there is a disadvantage that all of the three materials are contacted by a point contact, and thus the contact resistance is large.

Electrode for solid-state battery and manufacturing method therefor


A distinct cultivar of Chrysanthemum plant named xe2x80x98Iconxe2x80x99, characterized by its anemone-type inflorescences that are about 5.1 cm in diameter; attractive lavender purple-colored quilled ray florets and yellow green-colored disc florets; response time about 58 days; very freely flowering habit; dark green foliage; strong stems; and good postproduction longevity with inflorescences maintaining good substance and color for about two weeks in an interior environment.

The present invention relates to a new and distinct cultivar of Chrysanthemum plant, botanically known as Chrysanthemumxc3x97morifolium and hereinafter referred to by the name xe2x80x98Iconxe2x80x99.
The new Chrysanthemum is a product of a planned breeding program conducted by the Inventor in Salinas, Calif. and Alva, Fla. The objective of the breeding program is to create new cut Chrysanthemum cultivars having inflorescences with desirable colors and good form and substance.
The new Chrysanthemum originated from a cross made by the Inventor in September, 1995, in Salinas, Calif., of the Chrysanthemum cultivar Pilar, disclosed in U.S. Plant Pat. No. 10,331, as the female, or seed, parent with a proprietary Chrysanthemum seedling selection identified as 0384, not patented, as the male, or pollen parent.
The cultivar Icon was discovered and selected by the Inventor as a flowering plant within the progeny of the stated cross in a controlled environment in Alva, Fla., in November, 1996. The selection of this plant was based on its desirable inflorescence color and good form and substance.
Asexual reproduction of the new Chrysanthemum by terminal cuttings taken in a controlled environment in Alva, Fla. since February, 1997, has shown that the unique features of this new Chrysanthemum are stable and reproduced true to type in successive generations.
Plants of the cultivar Icon have not been observed under all possible environmental conditions. The phenotype may vary somewhat with variations in environment such as temperature, daylength and light intensity, without, however, any variance in genotype.
The following traits have been repeatedly observed and are determined to be the unique characteristics of xe2x80x98Iconxe2x80x99. These characteristics in combination distinguish xe2x80x98Iconxe2x80x99 as a new and distinct cultivar:
1. Upright cut Chrysanthemum that is usually grown as a natural spray.
2. Very freely flowering habit, about 12 inflorescences per flowering stem.
3. Anemone-type inflorescences that are about 5.1 cm in diameter.
4. Attractive lavender purple-colored quilled ray florets and yellow green-colored disc florets.
5. Response time about 58 days.
6. Dark green foliage.
7. Strong stems
8. Good postproduction longevity with inflorescences maintaining good substance and color for about two weeks in an interior environment.
Plants of the new Chrysanthemum are most similar to plants of the female parent, the cultivar Pilar. In side-by-side comparisons conducted by the Inventor in La Ceja, Antioquia, Colombia, South America, plants of the new Chrysanthemum had smaller inflorescences, shorter disc florets and more yellow green disc florets than plants of the cultivar Pilar.
Compared to plants of the male parent, the seedling selection identified as code number 0384, plants of the new Chrysanthemum have smaller inflorescences and shorter disc florets. In addition, plants of the male parent and the new Chrysanthemum differ in inflorescence form as plants of the male parent have daisy-type inflorescences.


There is provided a method, and a network node for implementing the method, in a cellular telecommunications network. The cellular telecommunications network includes a first and second base station, the first base station having a first interface for communicating with a first cellular core networking node and a second interface for wirelessly communicating with the second base station, the second base station being mobile and having a first interface for wirelessly communicating with the first base station and a second interface for communicating with a User Equipment (UE), wherein the UE's user traffic is transmitted between the first cellular core networking node and the UE via the first and second base stations.

Cellular telecommunications networks typically comprise a plurality of base stations which each have a coverage area for communicating with a plurality of User Equipment (UE). Traditionally, a base station is in a fixed geographical location by being installed at a particular site or by being installed on a building. These fixed base stations utilize local (typically wired) power and data communication interfaces to receive electrical power and to connect to a cellular core network (often known as a backhaul).
Modern cellular telecommunications networks also provide mobile base stations. These mobile base stations are of a non-fixed geographical location and may change their location either by self-powered movement or by being installed on a movable vehicle (e.g. a train or car). The mobile base station also has a coverage area for communicating with a plurality of UEs, which is achieved in the same manner as for fixed base stations. The mobile base station's backhaul may be provided by a wired interface, but in many circumstances such wired interfaces are unsuitable. Accordingly, mobile base stations typically have a wireless communications interface for carrying backhaul traffic. These wireless communications interfaces connect to a "donor" base station, and the donor base station's backhaul connection to the cellular core network includes both traffic for its own UEs and traffic for the mobile base station's UEs.
As the mobile base station may move relative to the fixed base stations, such that it may move out of a coverage area of its "serving" donor base station, cellular networking protocols define a handover process such that the mobile base station may disconnect from its serving donor base station and connect to a target donor base station with minimal disruption. This process mirrors the corresponding process for UE handover between a serving and target base station, and the mobile base station therefore includes a subset of UE functionality to support this process.

British Telecommunications Public Limited Company


An apparatus and method for communication using a wireless power are provided. The apparatus includes an amplifier configured to amplify an input signal based on a power supplied to the amplifier. The apparatus further includes a control unit configured to detect a change in an impedance of a target device, and to change the power based on the change in the impedance. The apparatus further includes a demodulation unit configured to receive a message from the target device, and to demodulate the message based on the changed power.

1. Field
The following description relates to an apparatus and method for communication using a wireless power.
2. Description of Related Art
Wireless power transmission may overcome problems, such as, for example, the inconvenience of wired power supplies, and limits to existing battery capacities, with respect to various electronic devices including, for example, electric vehicles, mobile devices, and/or the like. Wireless power transmission may use resonance characteristics of radio frequency (RF) elements. For example, a wireless power transmission system using resonance characteristics may include a source device configured to supply power, and a target device configured to receive the supplied power.

Apparatus and method for communication using wireless power


The present disclosure provides systems, methods, and computer-readable storage media having functionality to prove immutability of blockchains without accessing user data. A user may submit data for storage to a data management server and the data management server may generate one or more data records corresponding to the data at a database and one or more blocks at a blockchain, each block corresponding to of the data records. Block information associated with the generated blocks may be transmitted to a remote computing device for storage at a database. Prior to storing the block information, the remote computing device may sign the data using a private key or other cryptographic technique. To validate a block, raw block information may be retrieved from the blockchain and compared to the signed block information. If the signed block information matches the raw block information, the block may be determined to be valid (e.g., unchanged).

Blockchain technology has rapidly advanced in recent years and is being used across a wide range of industries and use cases. One property of blockchain technology that has contributed to its rapid adoption is immutability (e.g., the ability of the blockchain to remain unchanged, unaltered, and indelible), which increases a level of trust for the blockchain (e.g., a level of trust that the data recorded to or represented on the blockchain has not been altered). Blockchains may be public or private. Public blockchains are typically decentralized and may be considered permissionless because anyone can participate (e.g., read, write, etc.) in the blockchain network and there is no centralized control of the blockchain network. On the other hand, private blockchains are typically permissioned blockchains that place restrictions (e.g., permissions) on who can participate in the blockchain network and utilize a centralized control structure. While the immutability property is important for all blockchain technologies, it is especially important for blockchains operated under a centralized control structure, such as private blockchains. Because private blockchains utilized a centralized control, it may be easier (relative to decentralized public blockchains) to alter private blockchains. Regardless of whether a blockchain is public or private, the success of the technologies developed for and the services supported by the blockchain may be dependent on the ability of users to trust the blockchain. Thus, the ability to prove immutability can increase the trustworthiness of a blockchain and promote more widespread use of blockchain technologies.

Systems and methods for proving immutability of blockchains


Instructions in a computer system are executed in a plurality of parallel execution pipelines, a horizontal dependency check is carried out between instructions supplied to the parallel pipelines and in response to detecting horizontal dependency a control signal of a first or second type is generated depending on whether the dependency can be resolved by activating a by-pass or whether a temporary stall is required in one of the pipelines.

Existing machines are known to operate in superscalar mode in which two or more instructions are fetched from a program memory in the same machine cycle and then executed by different execution units. Commonly such instructions are located in adjacent memory positions and fetched from memory by a single fetch operation. Machines are also known to operate in very long instruction word (VLIW) mode. Such words may include a larger number of instructions retrieved from memory as a single word at a single word address. Instructions which may be included together in a single VLIW instruction are carefully controlled by software techniques to ensure that they avoid compatibility problems or data dependency problems between the instructions forming the word.
In computer systems a plurality of instructions may be executed in parallel and in pipelined operations. Data dependency may arise between two or more instructions which are executed in parallel or in a pipelined operation.
It is an object of the present invention to provide an improved computer system for detecting and resolving data dependence between instructions for parallel execution. Embodiments of the invention may provide a very high rate of instruction execution while avoiding a complex programming model for the software tool chain and assembly programmer.
In real-time operation the instruction scheduling should be deterministic, i.e., that the hardware does not do too much "intelligent" rescheduling of instructions which are uncontrollable (and non-deterministic) because this makes the programmer/software toolchain's task extremely difficult when the real-time nature of the application has to be respected.

Computer system that operates in VLIW and superscalar modes and has selectable dependency control


The present invention provides a method of manufacturing a printed circuit board. The method includes the steps of preparing an insulating substrate having a front surface and a back surface and a layer of metal foil formed on each of the front surface and the back surface; selectively forming a plating layer for forming a land on at least one of the metal foils; adjusting a thickness of the plating layer; and forming the metal foils into lines.

The present invention relates to a method of manufacturing a printed circuit board capable of adjustment of the aspect ratio of via holes and to a printed circuit board manufactured by the manufacturing method.
In recent years, the pitch between terminals of semiconductor chips such as BGA (ball grid array) packages has been reduced. There is a need to reduce the pitch between lines/lands provided on printed circuit boards for mounting of semiconductor chips as well as to reduce the diameter of via/through holes.
A land is a round or polygonal element of a circuit pattern provided on an insulating substrate for a printed circuit board. Terminals of a semiconductor chip are connected to lands, and via holes or through holes are formed in lands. Lines are portions of a circuit pattern other than lands.
When lines are formed with a reduced pitch, it is desirable to reduce the thickness of the metal conductor layer. This is because if the thickness of the metal conductor layer is thick when thin lines having a reduced line spacing (hereinafter referred to as fine lines) are formed by etching, etching cannot be suitably performed so as to avoid occurrence of metal residues (including an anchor) and failure to insulate lines from each other.
In a printed circuit board, openings for via holes or through holes are formed in an insulating substrate and plating is performed on the openings to form via holes or through holes. However, when plating is performed on the openings for via holes or through holes, it is performed on the entire insulating substrate. If plating is performed on the metal conductor layer in an area where lines will be formed, the entire metal layer becomes thick and it is difficult to reliably perform etching so that the lines have a predetermined line and a line spacing. This is a disadvantage to forming of fine lines.
On the other hand, as printed circuit boards become finer, the via hole diameter is reduced. In order to reduce the aspect ratio of via holes, it is necessary to reduce the thickness of the insulating layer (or dielectric layer). The aspect ratio is defined as the quotient of division of the depth of a via hole by the diameter of the via hole. If the insulating layer is a prepreg containing glass cloth, the thickness of the insulating layer is limited by the thickness of the glass cloth and cannot be reduced. In a case where via holes are filled with a plating, it is advantageous to reduce the diameter of the via holes from the viewpoint of the manufacturing time and the plating filling effect. If the diameter of the via holes is reduced, the aspect ratio of the via holes becomes higher and the plating liquid cannot flow smoothly to the bottom of each via hole. As a result, the plating layer at the bottoms of the via holes is made nonuniform in thickness and there is a possibility of occurrence of a crack in the plating layer at the reflow temperature at the time of mounting or in a thermal cycle test and, hence, a possibility of failure to establish interlayer connection.
In order to establish interlayer connection with reliability, it is necessary to increase the plating layer of in the via holes. However, the metal conductor layer in an area where lines are to be formed is also plated simultaneously with plating on the via holes to make the entire metal layer thick, as described above.
This is a disadvantage to forming of fine lines. As described above, it is difficult to reliably perform plating on via holes and to form fine lines.
The thickness of the plating layer formed by plating processing in the process of manufacturing a printed circuit board varies largely. Lines formed on the printed circuit board vary largely in thickness since the plating layer is superposed on metal foil. Also, at the time of etching for forming the lines, failure to completely performing etching may occur due to variation in thickness of the plating layer and there is a possibility of short circuit between the lines in some place.
In some cases, fine lines are formed in such a manner that metal foil on the insulating substrate is separated and plating is then performed to form the lines. However, there is a possibility of the bond strength of the plating layer to the insulating substrate becoming lower than that of the metal foil, i.e., failure to obtain the desired metal bond strength. There is a risk of the plating layer being separated from the insulating substrate when the completed printed circuit board is subject to various tests including a thermal shock test. If the plating layer is separated, the printed circuit board cannot be used as the product. In a case where plating processing is performed directly on an insulating member, there is a possibility of a catalyst such as palladium used in plating processing remaining as a residue to cause insulation degradation. If the catalyst remains as a residue between conductors, there is a risk of occurrence of failure phenomenon such as short circuit between pads when gold plating processing is performed and insulation degradation caused by an electrocorrosion test.
Japanese Patent Laid-Open Publication No. 50-41056 and Japanese Patent Laid-Open Publication No. 01-129494 disclose methods of manufacturing a printed circuit board by performing plating two times on through holes. It is possible to prevent occurrence of a crack in the plating layer by performing plating two times on through holes. Interlayer connections via the through holes can be established with reliability. However, plating is effected on metal foil as well as on through holes. Therefore, it is difficult to form fine lines even if interlayer connection can be established with reliability.
Japanese Patent Laid-Open Publication No. 03-175695 discloses a method of manufacturing a printed circuit board by performing plating on through holes after forming a circuit pattern. Plating is performed two times to ensure interlayer connection between the through holes. However, formed lines are also plated, so that it is difficult to form fine lines. Further, according to a drawing in Publication No. 03-175695, a plating layer is also formed on a side wall of the circuit pattern. Lines are made thicker by the thickness corresponding to the plating layer. The difficulty in forming fine lines is further increased thereby.
Japanese Patent Laid-Open Publication No. 04-62892 discloses a method of plating on through holes. This patent publication relates to a challenge to improve the reliability of connection via a through hole. The challenge is accomplished as described below. A circuit pattern is formed on an insulating substrate, and a mask is provided on the circuit pattern. An opening is formed in the insulating substrate by an operation from above the mask, and the mask on the periphery of the opening is also removed. Plating is performed on an inner opening surface to complete a through hole. According to FIG. 1 of Publication No. 04-62892, a portion on the periphery of the through hole is stepped. That is, the land at the through hole is stepped. The reliability of connection via the through hole in a thermal cycle test can be improved by stepping the land at the through hole. Publication No. 04-62892, however, includes no description of formation of fine lines and adjustment of the aspect ratio to a low value and differs from the present invention in terms of object.

Manufacturing method of printed circuit board

International Business Machines Corporation


An engine control system includes a cam phaser that introduces a cam phase angle θ between a camshaft intake lobe and an associated crankshaft. An engine control module communicates with the cam phaser to introduce the cam phase angle θ while an engine is being started. The cam phase angle θ is selected such that the camshaft intake lobe opens an intake valve during at least a portion of a compression stroke of a cylinder that is associated with the camshaft intake lobe.

The background description provided herein is for the purpose of generally presenting the context of the disclosure. Work of the presently named inventors, to the extent it is described in this background section, as well as aspects of the description that may not otherwise qualify as prior art at the time of filing, are neither expressly nor impliedly admitted as prior art against the present disclosure.
Spark-ignition direct-injection (SIDI) engines include one or more fuel injectors that inject fuel directly into associated engine cylinders. While the engine is running, a fuel pump supplies fuel to a fuel rail at high pressure, e.g. 3-15 megapascals (435-2176 pounds per square inch). The fuel pump can be mechanically driven by the engine. The fuel rail provides the pressurized fuel to the fuel injectors, which inject the fuel into the cylinders at times and pulse widths that are determined by an engine control module.
When the engine is being started it takes some time for the fuel pump to raise the pressure in the fuel rail. During that time the fuel pressure may be too low for the fuel injectors to provide enough fuel to start the engine quickly and reliably. A fuel flow rate through a fuel injector is based on a pressure difference across a fuel injector nozzle. When the engine is first being started, the engine is rotating slowly and the mechanically driven fuel pump may not develop fuel pressure that exceeds pressure in the cylinders. The fuel pressure therefore may not be able to overcome the cylinder pressures as is needed to deliver fuel to the cylinder via the fuel injectors.

Late intake valve closing to improve cold startability for spark-ignition direct-injection (SIDI) engines


Used deNOx or dedioxin catalytic converters are regenerated by washing the catalyst with a solution of surface-active substances in a liquid. At the same time, metal compounds are added for creating active centers.

The present invention relates to a process for the regeneration of used deNOx or dedioxin catalysts. Catalysts of this type are used in so-called deNOx or dedioxin installations for reducing and breaking down nitrogen oxides and/or in particular halogenated dioxins and furans in flue gases or other exhaust and off-gases.
The process known as selective catalytic reduction, or SCR for short, is one of the possible options for lowering or even substantially lowering the levels of nitrogen oxides NOx, i.e. a mixture of NO and NO2, formed for example during the combustion of fossil fuels in combustion plants. In the SCR process, the nitrogen oxides are converted into nitrogen and water using ammonia or substances which form ammonia under the system conditions as reducing agent and using a catalyst. Since the catalytic reactions proceed on the surface of the catalyst, a large specific surface area has to be provided through the use of correspondingly porous materials for the reaction. This requirement is met by the use of homogeneous ceramic catalysts, for example in honeycomb form. Most of a catalyst of this nature consists, for example, of the base material titanium dioxide TiO2 in which the active metal compounds, in particular V2O3, WO3, are homogeneously distributed. However, the catalyst may also be applied as a coating to a support, for example a metal sheet.
Under oxidizing conditions, identical or modified catalysts can also be used to lower the levels of emissions of organic products of incomplete combustion in off-gases from combustion plants, such as for example halogenated dioxins and furans. In this context, reference is had to the disclosure in international publication WO 91/04780.
There are transport processes upstream and downstream of the chemical reactions which take place on the catalyst surface. Following adsorption of the reaction partners on the internal surface of the catalyst, chemical combination between the reaction participants and the catalyst leads to a lowering of the activation energy which is absolutely imperative for the reaction to commence. A consequence is that the reaction is accelerated or the equilibrium is established.
If these active centers are blocked, for example by the accumulation of alkali metals and alkaline earth metals or their compounds which are contained in the fly ash, so that the activated NH3 adsorption required is partially impeded, the activity falls. In addition to this deterioration to the active areas of the catalyst surface through adsorbed catalyst toxins, the pores become blocked, for example, by calcium sulfate (CaSO4) and ammonium hydrogen sulfate (NH4HSO4) which are formed. Since the catalyst cannot be 100% selective with respect to a specific reaction, the catalyst also promotes some secondary reactions, including the conversion of SO2 to SO3, in an order of magnitude which is relevant. Although this reaction can be minimized by the composition of the catalyst, the fact remains that the small amount of SO3 is sufficient to react with the unreacted NH3, which is referred to as NH3 slippage, and H2O to form various salts, primarily to form ammonium hydrogen sulfate and ammonium sulfate (NH4)2SO4 or to combine with the fly ash.
These compounds form at temperatures at which condensation takes place when the temperature drops below the dew point of ammonium hydrogen sulfate. They may be deposited on the catalyst and in addition, together with adhesive particles, for example ash, fine dust, SiO2, Al2O3, may block the pores and thus lower the activity of the catalysts.
Therefore, the nature of the composition of the compounds which may be deposited on the catalyst is dependent on the composition of the fly ash, of the flue gas and of the operating temperature. They are generally alkali metal and alkaline earth metal compounds which are contained in the fly ash as oxides and, on account of their reaction with SO3, as sulfates and which are either deposited on the surface together with other compounds contained in the fly ash, such as for example SiO2 and Al2O3, and block the pores, or, on account of their electron donor properties, block the active centers and thus prevent the activated NH3 adsorption required.
The object of the present invention is to provide a method of regenerating a deNOx or deDioxin catalytic converter which overcomes the above-noted deficiencies and disadvantages of the prior art devices and methods of this general kind, and which provides for a process by means of which the number of the active centers available for the catalysis is increased as far as possible or as desired, namely, for example, up to the activity of the fresh catalyst or even beyond, in order in this way for the catalytic converter, i.e., the catalyst, to be fully or partially regenerated.
With the above and other objects in view there is provided, in accordance with the invention, a process for regenerating a used deNOx or dedioxin catalytic converters, which comprises washing a catalytic converter with a solution of surface-active substances in a liquid with a simultaneous addition of metal compounds creating active centers.
In other words, the above objects are achieved by the fact that the catalysts are washed with a solution of surface-active substances in a liquid, preferably water, with the simultaneous addition of metal compounds which create active centers.
As a result of this measure, deposited contaminants and chemisorbed compounds and ions are removed, old active centers are made available once more and additional active centers are created. In addition, in this way it isxe2x80x94quite surprisinglyxe2x80x94also possible to increase the activity compared to the fresh catalyst. The catalysts which have been treated in this way can be refitted into a deNOx, dedioxin or combined plant with their restored activity.
Washing in, for example, aqueous liquors is a complex operation in which numerous physical and chemical influences interact. This is understood as meaning both the removal of water-soluble deposits by water or by aqueous solutions of active washing substances and the detachment of water-insoluble deposits. In the process, it is possible to prevent redeposition of the insoluble fractions which have already been detached, for example by acoustic irradiation or by stabilizing the dispersed fractions. The water serves as a solvent for washing agents and for soluble compounds and as a transport medium for the dispersed fractions. The washing operation is initiated by the wetting and penetration of the substrate. This can be achieved quickly and completely if the high surface tension of the water is reduced substantially by surfactants as important washing agent components. The physical separation of the deposits from the substrate is based on the nonspecific adsorption of surfactants at various boundary surfaces which are present in the process. Substances with a low solubility are solubilized in molecularly dispersed form by surfactant micelles. The adsorption of washing agent constituents induces changes in the interfacial chemical properties and is consequently a precondition for good detachment.

Process for regenerating used deNOx or dedioxin catalytic converters

Integral Umwelt-und Anlangentechnik Gesellschaft m.b.H


A light source device includes: a power supply circuit, having a first voltage terminal and a second voltage terminal; a light source circuit, having a common terminal and multiple branch terminals, the common terminal being electrically connected to the first voltage terminal, the light source circuit including multiple light sources for emitting lights of different colors, first terminals of the light sources being electrically connected to the common terminal, and second terminals of the light sources being electrically connected to the branch terminals respectively; and a light adjusting circuit, including multiple electrical paths and being electrically connected to the second voltage terminal and the branch terminals via the electrical paths so that the light sources are connected in parallel between the first and second voltage terminals. A resistance of at least one of the electrical paths is adjustable to set branch currents respectively flowing through the light sources.

With the popularization of LED (Light emitting diode) products, there are increasing demands for adjustable color temperature of LED light sources. There are various requirements for adjustment of color temperature, some of which is that the color temperature is arbitrarily adjusted within a certain range, and some of which is that the color temperature is adjusted by setting multi-levels of color temperature within a range of color temperature. The conventional multi-level color temperature adjustment method is realized by re-designing a power supply circuit of a LED light source to add a control circuit in the power supply circuit. However, the added control circuit includes a MCU or a MOS transistor(s) with relatively high cost transistor(s), resulting in less design flexibility of the LED light source of multi-level color temperature adjustment and higher material cost. Therefore, how to realize an LED light source of multi-level color temperature adjustment at a relatively low cost is a technical problem needed to be solved.


Novel 4-substituted acetophenone anils and methods for preparing 1,3,5-tris(4'-hydroxyphenyl)benzenes from 4-substituted acetophenones such as 4-hydroxyacetophenones or, from substituted 4-hydroxyacetophenone-anils such as 4-hydroxyacetophenone-anil by reacting the 4-substituted acetophenone or corresponding anil with an aniline derivative.

1,3,5-tris(4'-hydroxyphenyl)benzene falls into the class of compounds known as trisphenyls. Trisphenyls have been recognized as useful intermediates in the preparation of more complex organic structures. For example, resins are readily prepared by a reacting trisphenyls with formaldehyde, acid anhydrides and more importantly with epichlorohydrin. Epoxide resins prepared from such compounds exhibit low shrinkage, extraordinary hardness, chemical inertness, outstanding mechanical strength, and a variety of beneficial features. See, for example, U.S. Pat. No. 4,394,496; and, the above-referenced application of Hilton.
1,3,5-tris(4'-hydroxyphenyl)benzene (THPB) molecules are particularly useful in their ability to stabilize polycarbonates. This is accomplished via a three site rigid D.sub.3h crosslink. THPB molecules may also be used as crosslinking agents in epoxy resins. See, for example, Chem. Abstracts., 66, 3004C.
THPB was reported in Beilstein, E II 6, 1115 (1921). The treatment of 4-methoxyacetophenone (4-MAP) with sulfuric acid produced 1,3,5-tris(4'-methoxyphenyl)benzene (4-MAP trimer or TMPB)(20% yield). This compound was demethylated with concentrated hydrochloric acid to yield THPB.
THPB was also reported in Chimia, 12, 143 (1958) and in Chimia, 13, 105 (1959) as formed by the trimerization of 4-haloacetophenone, where the halogen is either bromine or chlorine, in the presence of potassium pyrosulfate and sulfuric acid. This reaction results in 1,3,5-tris(4'-halophenyl)benzenes. These halogen-containing trimers were treated with sodium hydroxide and converted to THPB.
M. H. Karger and Y. Mazur, J. Org, Chem, 36, 540 (1971), reported that anisole and acetyl methanesulfonate, affords 4-MAP (46% yield) and TMPB (41% yield). Subsequent to anisole acetylation, trimerization is catalyzed by methanesulfonic acid.
R. E. Lyle, E. J. DeWitt, N. M. Nichols, and W. Cleland, J. Amer. Chem. Soc., 75, 5959 (1953), report the trimerization of substituted acetophenones, i.e., 4-MAP to TMPB (54% yield), by an alcoholic hydrogen chloride solution, after four months at room temperature.
G. P. Sharnin, I. E. Moisak, E. E. Gryazin, Zhurnal Prikladnoi Khimii, 43, 1642 (1970), report the trimerization of 4-MAP to TMPB (27% yield) using a mixture of potassium pyrosulfate and sulfuric acid. See also, A. F. Odel et al., J. Amer, Chem, Soc., 36, 81 (1913).
P. Milart and J. Cioslowski, Synthesis, p. 328-29 (1984) relate to the use of 4-alkoxyacetophenones to prepare 4-alkoxyacetophenone anils which are condensed to form 1,3,5-tris(4'-alkoxyphenyl)benzenes. However, the reference does not suggest using 4-hydroxyacetophenone (4-HAP) or substituted 4-hydroxyacetophenones to produce 4-hydroxyacetophenone-anil (4-HAP-anil) or substituted 4-hydroxyacetophenone-anils, which may be then condensed to form 1,3,5-tris(4'-hydroxyphenyl)benzene or substituted 1,3,5-tris(4'-hydroxyphenyl)benzenes. In fact, the reference teaches the conversion of 4-hydroxyacetophenone to 4-alkoxyacetophenone before converting to the corresponding anil and thereafter trimerizing the anil. This is consistent with earlier teachings which describe unsuccessful efforts to trimerize hydroxyacetophenone. See, for example, G. P. Sharnin et al., supra, see also R. E. Lyle et al., supra. Its is also consistent with prior teachings that 1,3,5-tris(4'-hydroxyphenyl)benzene is produced by hydrolyzing the corresponding 4'-alkoxy substituted compound which is prepared using 4-alkoxyacetophenone. See, for example, Beilstein, E II 6, 1115 (1921). Moreover, P. Milart et al. fail to teach or suggest a one step process for trimerizing 4-HAP.
U.S. Pat. No. 3,458,473, issued Jul. 29, 1969 to Starnes et al. is directed to the preparation of various hindered trisphenyls prepared by the cyclotrimerization of an acetylphenol precursor.
U.S. Pat. No. 3,644,538 issued Feb. 22, 1972 to W. H. Starnes, discloses that both 3'-alkyl- and 3',5'-dialkyl-4'-hydroxyacetophenones can be trimerized to the corresponding triarylbenzenes with anhydrous HCl and triethyl orthoformate and ethanol. Trimerization of unsubstituted 4-HAP is not suggested or disclosed. Starnes also fails to teach or suggest trimerizing substituted or unsubstituted 4-HAP-anil.
German Patent 258,929 to Zimmerman et al., issued Aug. 10, 1988 is directed to methods for the production of 1,3,5-tris(triarylbenzene) compounds. These compounds are reacted by combining 2,4,6-triaryl pyrylium salts with carboxylic acid anhydride in the presence of a basic condensing agent. The reaction of Zimmerman utilizes triaryl pyrylium carboxylic anhydride.
Elmorsy et al., "The Direct Production of Tri- and Hexa-Substituted Benzenes from Ketones Under Mild Conditions," Tetrahedron Letters, Vol. 32, No. 33, pp. 4175-4176 (1991) report the treatment of aryl benzenes with tetrachlorosilane in ethanol to yield 1,3,5-triarylbenzenes. However, Elmorsy et al. fail to teach or suggest trimerizing 4-hydroxyacetophenone, a 4-hydroxyacetophenone derivative, or a 4-substituted-oxyacetophenone by contacting such a compound with a halosilane, as in the present invention. Indeed, it is believed that trimerizing a hydroxyacetophenone such as 4-hydroxyacetophenone or a 4-hydroxyacetophenone derivative, is not disclosed or suggested by Elmorsy et al. because of the belief that the hydroxy group would interfere with the reaction, for example, react with the tetrachlorosilane. See, e.g,, Sharin et al., supra, which illustrate why it was believed, before now, that direct trimerization of hydroxyacetophenone was not feasible. Accordingly, Elmorsy et al. fail to teach or suggest the present invention.

Process for the prepartion of 1,3,5-tris(4'-hydroxyaryl)benzene


A method of determining an adjusted color to be used for computing colorants for printing on a specified substrate, comprising: specifying an apparent color; estimating diffuse reflection from an outside surface of colorants when printed on the specified substrate; and adjusting the specified color for the effects of the estimated diffuse reflection to determine a color to be used for computing the colorants.

U.S. Pat. No. 4,035,214 to Shuppert, et al. describes a printing system in which a release blanket is coated with a transparent film. A multi-color image is printed onto the release film and film, together with the images is transferred, image side down, onto a can. The release thus acts as a protective coating for the image. A film of adhesive is applied either to the image or to the can so that the film will transfer from the release blanket. The images are preferably printed by flexographic, letter press or other systems utilizing raised ink applicators.
U.S. Pat. No. 4,949,103 describes a direct electrostatic printing apparatus in which charged toner is applied by a print-head structure to a film coated with an adhesive to produce a "reverse reading image." The film with the image thereon are adhered to a final surface, with the image facing the final surface by the adhesive. The printed images may be alphanumeric text (for use in address labels) or a picture (for use in an identification badge).
It is well known in the art to print various types of images on paper or some other surface and then to laminate the printed image using either a single sheet of plastic material attached to the printed surface or by laminating a printed sheet between two sheets of plastic. In general, the plastic sheets are coated on one side with an adhesive material which adheres to the image. In some prior art systems the adhesive is not sticky at room temperature and the lamination process is carried out which subjecting the plastic sheet and the printed substrate to heat and pressure.
However, when thick plastic material is used, the image appears to be behind the front surface, which is not esthetically pleasing. On the other hand, when utilizing thinner plastic sheets, the structure or substructure of the printed sheet transfers to the front surface, reducing the quality of the resulting product.
These problems form a limitation on the use of such lamination in producing digital photographic prints (referred to herein as "digital photo-finishing"), i.e., the limitations on the surface finish available. In particular, laminated images, even when produced with very smooth plastic film, tend to have, on them, a pattern or other variations, unless they are thick. These variations on the surface may be due to the surface of the film following the surface texture of the paper to which it is bonded, or other causes.
WO 99/19733 and PCT applications PCT/IL00/0084, entitled "Printing on Plastic" and PCT/IL99/00510, entitled "Substrate Coating for Improved Toner Transfer and Adhesion", the disclosures of all of which are incorporated herein by reference, describe coating of plastic sheets with layer(s) which adhere to the sheets and to which liquid toners can be transferred and bonded. In general, these coatings comprise thermoplastic material.
It is known in the art to coat paper or other materials with thin layers of plastic. One method of coating, called casting, comprises extruding a thin layer of the plastic of the paper on the paper and then passing the coated paper through a pair of rollers, while the plastic is still hot. The plastic replicates the finish of the roller that touches it. One use for such coated paper is as a substrate for various printing processes.

Hewlett-Packard Development Company, L.P.


A conventional game program involves the risk of the player losing interest in the game. A computer-readable program product which enables a beginner player to sufficiently enjoy interest in a game is provided. A character to be subjected to ring customization is selected. If a technical ring is selected as a ring type after the subject has been determined to be change of a ring type, an ally character makes attacks only when the player has performed operation input with a rotary bar situated in any of timing areas. When a gamble ring is selected as the ring type, only one timing area is available. In a case where a practice ring is selected as the ring type, determination of operation input timing is continuously made in all the remaining timing areas even when the operation input has been performed with the rotary bar not situated in any timing area.

1. Field of the Invention
The present invention relates to a computer-readable program product for implementing a game in which a plurality of characters are displayed on a display screen and a player is caused to select action patterns of the characters, to thus cause the game to proceed. The present invention also relates to a gaming machine constituted of a computer or specific hardware for executing the game.
2. Description of the Prior Art
Various games have already been provided, wherein a virtual world of the game appears on a screen of a computer or a TV (television) set. In the virtual world, a player controls characters in the game by means of entering commands by way of operation means, such as a controller or the like, thereby causing a predetermined story to proceed. These games are generally referred to as "RPGs" (Role Playing Games).
A commonly known RPG is as follows (see, e.g., Japanese Unexamined Patent Publication No. 2002-200354). Specifically, a combat situation is embedded in the RPG, wherein a character (hereinafter called an "ally character") controlled by the player battles against enemy characters controlled by a computer. The ally character acquires an experience value and virtual currency by means of defeating the enemy characters in battle. Thus, a story is caused to proceed while the level of the ally character is being raised.
In the combat situation in the RPG described in Japanese Unexamined Patent Publication No. 2002-200354, the pattern of an attack delivered by the ally character is determined in accordance with operation of buttons performed by the player after the player has selected an "attack" command. Specifically, when the player has selected an "attack" command, a variable display zone called a judgment ring appears on the screen, and a variable zone called a rotary bar rotationally moves around the variable display zone. A plurality of effective zones are formed in the variable display zone. In a case where the player has operated the buttons at a good timing when the rotary bar has passed through one of the effective zones, the ally character commences an attack. The pattern and efficacy of the attack are determined in accordance with an ability corresponding to the level of the ally character and items on hand (weapons, magic (spell), or the like). However, when the buttons are operated at a timing at which the rotary bar is in a special effective zone of the effective zone, the degree of effectiveness of the attack; that is, the efficacy of the attack, is enhanced. In contrast, in a case where the player has operated the buttons when the rotary bar is not situated in the effective zone of the variable display zone, the ally character does not make an attack. Moreover, the number of attacks corresponds to the number of times the player has operated the buttons at a good timing when the rotary bar is situated in any one of the effective zones.
In a case where the buttons are operated when the rotary bar is not situated in any one of the effective zones, the operation input is terminated at that point in time even when the effective zones—which the rotary bar has not yet passed through—still remain. When the player has failed to perform a button operation input within the effective zone which the rotary bar is to pass through first, the ally character does not make an attack.
As mentioned above, in the conventional game program, even when the buttons are operated during the course of the rotary bar passing through a plurality of the effective zones and the rotary bar is not situated in any one of the effective zones, the ally character carries out attacks by the amount corresponding to the number of times the player operates the buttons at a good timing. Consequently, some of the conventional game program involve the risk of the player losing interest in the game, because the player cannot find any amusement in operating the buttons.
Moreover, in the above-described conventional game program, the player can operate the buttons for determining the number of attacks for a plurality of respective effective zones, thereby relieving the player's strain on operation of the buttons. For this reason, in the conventional gaming machine, the player cannot find amusement in operating the buttons. Consequently, the gaming machine involves the risk of the player losing interest in the game.
In the conventional game program, it is not easy for a beginner player to operate the buttons when the rotary bar passes through any one of the effective zones. When the beginner player plays a game, the number of attacks made by the ally character is reduced, thereby posing difficulty in causing the story to proceed and preventing the player from sufficiently enjoying interest in the game.

This patent is subject to a terminal disclaimer.

Program for implementing a game and gaming machine


Systems and methods to determine random coincidence rates include determination of a detector rate for each of a plurality of detectors of a positron emission tomography scanner based on a frame of positron emission tomography data, determination of a sensitivity for each detector crystal of the plurality of detectors, based on the detector rate of the detector including the detector crystal, determination of a singles rate for each detector crystal based on the detector rate of the detector including the detector crystal and the determined sensitivity of the detector crystal, estimation of a mean random coincidence rate for each of a plurality of pairs of the detector crystals based on the singles rate of each detector crystal of each of the plurality of pairs of the detector crystals, correction of the acquired frame of positron emission tomography data based on the estimated mean random coincidence rates, and reconstruction of a positron emission tomography image based on the corrected frame of positron emission tomography data.

According to conventional positron-emission-tomography (PET) imaging, a radiopharmaceutical tracer is introduced into a patient body typically via radial arterial injection. Radioactive decay of the tracer generates positrons which eventually encounter electrons and are annihilated thereby. Annihilation produces two photons which travel in approximately opposite directions. A ring of detectors surrounding the body detects the emitted photons and reconstructs PET images based thereon.
A "coincidence" is identified when two detectors disposed on opposite sides of the body detect the arrival of two photons within a particular coincidence time window. Because the two "coincident" photons travel in approximately opposite directions, the locations of the two detectors determine a Line-of-Response (LOR) along which an annihilation event may have occurred.
A "true" coincidence represents the detection of two coincident photons which arose from a single annihilation event located on a LOR between the two detectors. A "random" coincidence represents two coincident photons which did not arise from the same annihilation event. A "scatter" coincidence is a type of true coincidence in which two coincident photons originated from the same annihilation event but the annihilation event was not located along the LOR of the two detectors because one or both of the photons interacted and scattered within the body or media. Since only the true unscattered coincidences indicate locations of annihilation events, random coincidences and scatter coincidences should be subtracted from or otherwise used to correct acquired PET data reconstruction of a PET image.
Conventional PET scanners detect all coincidences without regard to whether the coincidences are true, random or scatter coincidences. Software and/or hardware-based approaches can be used to estimate random coincidences and to subtract the random coincidences from the detected coincidences. For example, one current hardware-based approach involves delaying one of the detected singles events. Undelayed logic detects all coincidences along all LORs as described above. Additional delayed logic delays one input channel by, for example, a few tens of nanoseconds (e.g., 5× the coincidence window) and then performs coincidence detection. As a result, the delayed logic does not detect any actually-true coincidences as coincident. To correct for random coincidences, and along each LOR, the coincidences detected by the delayed logic are subtracted from those detected by the undelayed logic.
Direct subtraction of the random coincidences increases statistical noise due to the statistical variation of each detector pair. Moreover, duplication of the coincidence detection by the coincidence detection components limits the count rate capability of a scanner, particularly in high-activity PET studies. Systems are desired to efficiently correct acquired coincidences data for the presence of random coincidences.

Determination of crystal singles rates to estimate mean random coincidence rate


To provide a fixing device for a steering device capable of tilt adjustment and having a very simple cam structure.

1. Field of the Invention
The present invention relates to a fixing device for a steering device capable of tilt adjustment and having a very simple cam structure.
2. Description of the Related Art
Among tilt telescopic steering devices in which the height of a steering wheel and the orientation of the axis can be changed as appropriate to the body build and driving attitude of the driver, steering devices provided with a tilt adjustment mechanism whose structure includes an operation lever fitted to a lock bolt that passes through side plates of a fixed bracket installed on the car body and side plates of a movable bracket in which the steering column is installed, and by rotating the operation lever the side plates of the fixed bracket and movable bracket press against each other to lock the tilt are known.
Among this type of device are devices that include two cam members: one cam member fitted to the fixed bracket side, and the other cam member fitted to the operation lever for tilt adjustment, such as Japanese Patent Application Laid-open No. 2002-59851. This has a structure in which as a result of rotation of the operation lever, the cam surface of one cam member acts on the cam follower surface of the other cam member, so that the other cam member presses against the side part of the fixed bracket and the fixed bracket fixes the movable bracket. The cam member fitted to the operation lever is screwed to the threaded part of the tilt bolt. When the operation lever is rotated one cam member rotates, and acts on the cam follower surface of the other cam member fixed to the fixed bracket side part, so that the movable bracket side part presses against the fixed bracket side part.
This conventional cam configuration is shown in more detail in FIGS. 13A through 13C. The configuration includes a main drive cam 50 having a cam follower surface 51 (plateau 51a, slanting part 51b, valley 51c) and a driven cam 60 having a plurality of projections 61 that move up and down along the cam follower surface 51. The main drive cam 50 and driven cam 60 are fitted together. When the main drive cam 50 rotates relative to the driven cam 60, the locked state is when the projections 61 of the driven cam 60 rise onto the plateau 51a of the cam follower surface 51 of the main drive cam 50, and the released state is when the projections 61 of the driven cam 60 move from the plateau 51a of the cam follower surface 51 to the valleys 51c.

Abstract: The fixing device includes a fixed bracket; a movable bracket; a lock bolt that passes through the fixed bracket and the movable bracket; a main drive cam, through which the lock bolt passes, and that is fitted to rotate freely; and a driven cam, through which the lock bolt passes, and that is fitted unrotatably. Cam operation parts are formed on the main drive cam and the driven cam respectively. Cam operation parts on both sides are always in contact with the other side while the main drive cam and the driven cam separate and come close in the axial direction, and the diameter of the driven cam is larger than the diameter of the main drive cam.