An X-cut or near Y-cut monocrystalline wafer of electro-optical material forms a substrate of a digital optical switch. Waveguides of LiNbO.sub.3 are diffused into the wafer to form ribbon waveguides such that light propagates in a direction near the Z axis of the wafer. These waveguides form a Y-shape, refractive indices of which are controlled by electrodes formed closely adjacent thereto.

1) Field of the Invention
The invention relates to an integrated optical circuit using thin film optical waveguides for modulating or switching a light signal by an electro-optic effect and, more particularly, to a X- or almost Y-cut, near Z propagating digital optical switch.
2) Description of Related Art
The growing utilization of optical fibers in communications, sensors and other applications has made optical switching a subject of great interest. Of particular interest is the switching of signals in optical form without electro-optical transformations. Attempts have been made in a number of directions, such as by using mechanical switches and integrated optical semiconductor switches made of semiconductor material such as GaAs and InP.
Many of these switches use an electro-optical effect for the switching function. A quadratic or Kerr electro-optical effect is present in all substances and refers to a change in refractive index .DELTA.n proportional to a square of the applied electric field E. Much larger index changes can be realized in single crystals that exhibit a linear or Pockel's electro-optic effect. In this case the change of refractive index .DELTA.n is directly proportional to the applied electric field E. The effect is present only in noncentrosymmetric single crystals, and the induced index change depends upon the orientation of the electric field E and the polarization of the light beam. Well known linear electro-optic materials include potassium dihydrogen phosphate (KDP) and its deuterated isomorph (DKDP or KD*P), lithium niobate (LiNbO.sub.3) and lithium tantalate (LiTaO.sub.3), and semi-conductors such as gallium arsenide (GaAs) and cadmium telluride (CdTe).
To date, one of the most mature technologies for fabrication of integrated optics circuits (IOC's) is Ti:LiNbO.sub.3, which involves a ribbon guide formed by diffusing titanium into lithium niobate. The lithium niobate guide is formed in a monocrystalline wafer upon which longitudinal electrodes are placed adjacent the ribbon guide. The electrodes modify the local value of the refractive index when a potential difference is applied.
In this technology, relatively low loss channels with good spot-size matching to single mode fibers can be conveniently fabricated. However, there is one important problem; the LiNbO.sub.3 components are inherently polarization dependent, i.e., they require excitation of a specific linear polarization ("ordinary" transverse electric mode TE or "extraordinary" transverse magnetic mode TM, for a Z-cut example), dependent on crystal orientation, whereas the state of polarization (SOP) at the output of a single mode fiber will exhibit a random behavior. This problem is usually solved by using techniques adapted to accommodate the inherent polarization dependency of the LiNbO.sub.3 components, which serves to complicate the systems.
There are a number of concepts for polarization independent switches and modulators. See, e.g., M. Kondo et al, "Low Drive Voltage and Low Loss Polarization Independent LiNbO.sub.3 Optical Waveguide Switches," Electron. Lett., Vol 23, (1987), pp 1167-1169; R. C. Alferness, "Polarization Independent Optical Directional Coupler Switch Using Weighted Coupling," Appl. Phys. Lett., Vol 35, (1979), pp 748-750; O. G. Ramer et al., "Polarization Independent optical Switch with Multiple Sections of .DELTA..beta. Reversal and a Gaussian Taper Function, " IEEE Journ. Quantum Electron. Vol QE-18 (1982), pp 1772-1779; L. McCaughan, "Low Loss Polarization Independent Electrooptical Switching at .lambda.=1.3 .mu.m," IEEE Journ. Lightwave Techn., Vol LT-2, (1984), pp 51-55; Y. Bourbon et al., "Polarization-Independent Modulator with Ti:LiNbO.sub.3 Strip Waveguides," Electron. Lett. Vol 20, (1984), pp 496-497; N. Tsukada et al., "Polarization-Insensitive Integrated-Optical Switches: A New Approach," IEEE Journ. Quantum Electron., Vol QE-17, (1981), pp 959-964; J. E. Watson, "A Low-Voltage Polarization Independent Guided Wave Direction-Coupler switch in Lithium Niobate," SPIE Vol 835, Integrated Optical Circuit Engineering V, (1987), pp 132-135; J. E. Watson et al., "A polarization Independent 1.times.16 Guided-Wave Optical Switch Integrated on Lithium Niobate," Journ. Lightwave Techn., Vol LT-4, (1986), pp 1717-1721; W. K. Burns et al., "Interferometric Waveguide Modulator with Polarization-Independent Operation," Applied Physics Letters, Vol 3, (1978), p 944; P. Granestrand et al., "Polarization Independent Optical Switches," Fourth European Conference on Integrated Optics (ECIO '87) pp 36-39; P. Granestrand et al., "Polarization Independent Switch and Polarization Splitter Employing .DELTA..beta. and .DELTA..beta. Modulation, " Electron. Lett. 1988, 1142-1143; J. L. Nightingale et el, "Low-Voltage Polarization Independent Optical Switch in Ti-indiffused Lithium Niobate," Techn. Digest of Integrated and Guided Wave Optics Conf. (IGWO '89), paper MAA3, pp 10-13; K. Takizawa et al., "Polarization-Independent and Optical Damage-insensitive LiNbO.sub.3 Interferometric Waveguide Modulator," Japanese Journal of Applied Physics, Vol 27, (1988), pp L696-L698; Y. Silberberg et al., "Digital Optical Switch," Techn. Digest OFC 1988, paper THA3; H. F. Taylor, "Polarization Independent Guided wave Optical Modulators and Switches," IEEE Journ. Lightwave Techn., Vol LT 3 (1985), pp 1277-1280; T. Pohlmann et al., "Polarization independent switches on LiNbO.sub.3," Proceedings of the Topical Meeting on Integrated Photonics Research, Hilton Head, SC, 1990, pp 38-39. The first experimental results on polarization independent switches were reported by Alferness in 1979.
An approach to making polarization independent switches is to utilize crystal orientations where the conditions are similar for the two polarizations. This means that the electro-optically induced perturbations are equal and that the TE and TM modes have approximately the same coupling lengths. The "isotropic" orientations with the Z-axis in the propagation direction are examples of such orientations. Here both polarizations see the ordinary refractive index and therefore the index perturbations due to titanium (Ti) concentration are equal for the two polarizations. This means that the coupling lengths are approximately equal. The electro-optically induced index perturbations for the two polarizations are caused by the electro-optic r-coefficients r.sub.12 and r.sub.22 (contracted index notation). They have equal magnitudes but opposite signs and the index 2 which is common in r.sub.12 and r.sub.22 implies that these perturbations correspond to external electrical fields along the Y-axis. One complicating factor in this context is the fact that the two polarizations are almost synchronous and that there is one electro-optic coefficient (of the same magnitude as r.sub.12 and r.sub.22) which performs a coupling between the two polarizations. This r-coefficient is r.sub.61, indexed 1, 2, 1 in non-contracted notation. The index 6 (1,2 non-contracted) corresponds to coupling between electrical fields along the X and Y axis directions (the notation format is explained below). The coupling is induced by an external electrical field along the X axis as indicated by the second listed index.
To get good performance in a switch in this orientation, this TE-TM conversion must be avoided. However, since the unwanted (TE.fwdarw.TM) and the wanted (.DELTA.n.sub.TM, .DELTA.n.sub.TE) perturbations corresponds to different components of the external field, it is possible to avoid this coupling by proper design of the component.
One advantage with Z-propagating concepts is that the same index (ordinary) is seen by the two different polarizations. Therefore, there will be no bandwidth degradation due to pulse broadening which occurs when different indices are seen by the two polarizations.
This index difference (which implies different velocities for the portions of a pulse in the TE and TM polarizations, respectively) will appear when crystal orientations with the Z axis perpendicular to the propagation direction is used and will limit the permitted bitrate to approximately 10 Gbit/s per channel for a 10 cm long chip in these orientations. However, if the information is stacked by some other principle, such as wavelength division multiplexing or by coherent techniques, information band-widths in the THz-range can be switched even with switches in these orientations.
Another advantage with the Z-propagating concepts is that both polarizations will have approximately the same transfer function (in general, for other concepts the polarization independency means that it is possible to put the switch in two polarization independent switch states but without independency in "intermediate" points) This is most important, e.g., when linear (small signal) modulation applications are considered.
All of the switch types mentioned above are so called interferometric switches; they are based on constructive and destructive interference of modes, therefore they all have oscillating transfer functions. There is, however, another possibility, and that is to use devices which are based on mode sorting instead of interferometry. See, W. K. Burns et al., "Mode Conversion in Planar Dielectric separating Waveguides", IEEE Journ. of Quantum Electron., Vol QF-11 (1975), pp 32-35, and Y. Silberberg et al., "Digital optical Switch," Tech. Digest OFC, 1988, paper THA3.
FIG. 1 shows a mode sorting switch 100. It is a 2.times.2 digital optical switch and consists of an asymmetrical Y-branch waveguide 101 at one side and a symmetrical Y-branch waveguide 102 at the other. The latter branch 102 can be made asymmetric by applying an electrical field via the electrodes 103 and 104. An asymmetric Y-branch 101 performs mode sorting, given that the transformation is "adiabatic" (sufficiently slow). The mode sorting here means that the channel mode in the input waveguide with highest effective index gradually transforms along the branch to the first order local normal mode (fundamental mode) of the two mode region where the channels are near and influence each other (there is no power transfer between the local normal modes; at large separation the first order mode has the shape of the channel mode). In the same way, the mode in the other channel transforms to the second order mode (i.e., first higher order mode). Thus, the signal in the wide channel transforms to the first order mode in the middle region and the signal in the narrower channel transforms to the second order mode of the middle region. If the other half 102 of the switch 100 is also asymmetric (i.e., an odd index perturbation is induced) by applying a voltage to the electrodes 103 and 104 in a similar way, the first order mode in the middle region (corresponding to the wide input channel) transforms to the output channel with highest index and conversely the second order mode transforms to the output channel with lowest index.
Since the output Y-branch asymmetry can be electro-optically altered the device works as a 2.times.2 switch provided that the index perturbation is large enough and the transformation is adiabatic. If zero voltage is applied to the electrodes, a 3 dB splitting will occur for both signals.
In FIG. 2 a transfer function for a X-cut digital optical switch for TE (solid line) and TM (dashed line) is shown (note the difference). As can be seen, the transfer function does not have the oscillatory behavior of interferometric switches and the switch will operate independent of polarization assuming that the magnitude of the drive voltage is high enough.
A significant advantage of the digital optical switch is its superb stability performance. The instabilities caused by DC drift and temperature variations appear as variations of the "effective" applied voltage, and the digital response with its small transfer function slope attenuates the induced switch state perturbation if an operating point with sufficiently high voltage magnitude is chosen.
Another advantage of the digital optical switch is the power splitting achieved at zero voltage. This is especially important when broadcasting operation is requested, e.g., in some switch matrix applications.
FIG. 3 shows a switch matrix structure in which the digital switch is a very attractive choice of switch element. See, R. A. Spanke, "Architectures for Large Non-blocking Optical Space Switches," IEEE Journ. of Quantum Electron., Vol QE-22 (1986), pp 964-967. Here 1.times.2 switches are needed, which means that the switch described above can be simplified to the 1.times.2 digital optical switch structure of FIG. 4 (X-cut example).
In the switch shown in FIG. 4, the signal in the incoming single mode channel 131 excites the first order mode of the two mode region 132, and this mode is transformed (as described above) to the output with highest effective index by application of an appropriate electric field via three electrodes 133, 134, 135 (i.e., 133 and 135 being grounded and 134 being at positive or negative potential with sufficient magnitude, or 134 being grounded and 133 and 135 being at positive or negative potential with sufficient magnitude). The structure of FIG. 3 has good cross-talk properties because a signal must go wrong in two switches before it reaches an unwanted output if the matrix is properly set (this is not true in passive splitting-active combining operation as further explained infra). See, the Spanke article.
The good cross-talk performance of the structure relaxes the cross-talk requirement on the individual switch elements, which is advantageous because it is probably more difficult to achieve extremely low cross-talk with the digital optical switch than with an electronically adjustable directional coupler such as shown in the Granestrand et al. article.
Another important feature of a switch matrix due to FIG. 3 with digital optical switches as switch elements is the possibility to conveniently implement broadcast functions where the signal from one input is distributed to several outputs. In this case some (or all) of the switches in the first half of the matrix is set to 3 dB power splitting which in the case of the digital optical switch is reached at zero voltage.
Each of the couplers or switches described above suffer drawbacks such as complex fabrication, elaborate implementation and/or control, limited bitrate capacity and undue sensitivity to the environment.


A semiconductor device of the present invention has a first conductive layer, a second conductive layer, an insulating layer which is formed between the first conductive layer and the second conductive layer and which has a contact hole, and a third conductive layer which is connected to the first conductive layer and the second conductive layer and of which at least a part of an end portion is formed inside the contact hole. Near a contact hole where the second conductive layer is connected to the third conductive layer, the third conductive layer does not overlap with the second conductive layer with the first insulating layer interposed therebetween and an end portion of the third conductive layer is not formed over the first insulating layer. This allows suppression of depression and projection of the third conductive layer.

1. Field of Invention
The present invention relates to a semiconductor device in which wires are stacked.
2. Description of the Related Art
In recent years, the degree of integration of semiconductor elements increases for wiring substrates and semiconductor devices, and integration progresses further. This makes wire width narrower and increases the number of wires. Moreover, downsizing of an area of semiconductor devices is attempted by making the wiring into multi-wiring.
As a typical example of a semiconductor element that forms a semiconductor device, a thin film transistor is described with reference to FIGS. 16A to 16C. FIG. 16A is a top view of a thin film transistor and FIGS. 16B and 16C are cross-sectional views taken along lines A-B and C-D of FIG. 16A, respectively. In FIG. 16A, a substrate 101, insulating layers 102, 107, and 115, and a gate insulating layer 104 which are shown in FIGS. 16B and 16C are omitted.
As shown in FIG. 16B, in a thin film transistor 110, a semiconductor region 103 and a gate electrode 105 are insulated by the gate insulating layer 104 and the gate electrode 105 and wires 108 and 109 are insulated by the insulating layer 107. Moreover, a conductive layer 106 is formed at the same time as the gate electrode 105 over the gate insulating layer 104, and the conductive layer 106 and the wire 109 are insulated by the insulating layer 107 interposed therebetween.
Moreover, the wire 108 is connected to one of a source region and a drain region of the semiconductor region 103 through a contact hole 111 formed in the insulating layer 107. Moreover, the wire 109 is connected to the other of the source region and the drain region of the semiconductor region 103 through a contact hole 112 formed in the insulating layer 107. The wire 109 is connected to the conductive layer 106 through a contact hole 113 formed in the insulating layer 107.
Such a structure makes it possible to stack a semiconductor region, a gate electrode, a wire, and the like in a thin film transistor.
As shown in FIG. 16A, each of the wires 108 and 109 is formed so as to have a cross sectional area that is larger than an area of each of the contact holes 111 to 113. In the case of forming the insulating layer 107 by a thin-film formation method such as a CVD method or a sputtering method, a surface of the insulating layer 107 is not flattened and is affected by depression and projection of a layer serving as a base for the insulating layer 107. Thus, depression and projection increase on the surfaces of the wires 108 and 109 and the insulating layer 107, as shown in FIG. 16B.
In a liquid crystal display device, if an element for driving a pixel electrode is formed by using such a thin film transistor, depression and projection of the insulating layer cause an orientation film not to be rubbed uniformly. This causes problems in that orientation of liquid crystals is distorted and accordingly image quality is degraded. Therefore, the insulating layer 115 which has been flattened is formed over the insulating layer 107 and the wires 108 and 109.
As a method for forming the insulating layer 115, a method in which an insulating film formed by a CVD method or a PVD method is flattened by polishing by CMP or the like or a method in which an insulating layer having superior flatness is formed by a coating method is given.
In the case of forming the insulating layer 115 having superior flatness by a coating method, the insulating layer 115 is formed in such a way that a composition that forms the insulating layer is applied by a coating method such as a spin coating method or a slit coating method and then the composition is baked. However, it is necessary to form the insulating layer 115 thick in a depression portion of the insulating layer 107 and the wires 108 and 109 that serve as a base for the insulating layer 115 though the insulating layer 115 may be thin in a projection portion thereof. As a result, as compared with the case of forming an insulating layer over an insulating layer and wires that have little depression and projection difference, the amount of compositions that form the insulating layer 115 increases and material cost increases, as well as the productivity is reduced. Moreover, in the case where the insulating layer 115 is formed with a photosensitive material and light-exposed in later steps, a problem occurs that a light-exposure time extends if the film is thick.
Further, as shown in FIG. 16B, depression and projection difference in a region 116 of the wire 109 where the thick conductive layer 106 overlaps with the insulating layer 107 is larger than in a region 117 of the wire 109 where the insulating layer 107 does not overlap with the conductive layer 106. Thus, uniform application of the composition is interrupted in the regions 116 and 117, resulting in that the film thickness of the insulating layer 115 decreases in the region 116. Accordingly, a problem occurs that a wire formed over the insulating layer 115 is easily short-circuited with the wire 109.

Semiconductor Energy Laboratory Co., Ltd.


A piston-cylinder unit for a mono-tube damper, includes: a cylinder having an inner wall and a closure fixed axially with respect to the cylinder by positive engagement with one or more impressions of the inner wall, the closure having a plug sealed with respect to the inner wall, the plug having a reinforced portion and being axially moveable with respect to the cylinder when a temperature threshold is exceeded. Axial movement caused by the temperature threshold being exceeded causes deformation of the plug with respect to the cylinder so as to: (i) form a channel between the plug and the inner wall of the cylinder, through which pressure, caused by the increased temperature, can escape the piston-cylinder unit, and (ii) permit axial movement of the closure only to a point at which the reinforced portion of the plug comes into contact with the one or more impressions in the cylinder.

1. Field of the Invention
The invention is directed to a cylinder piston device and more particularly to a mono-tube vibration damper, including a cylinder having an inner wall bounding a work space filled with a working medium, such as hydraulic fluid. A closure is fixed axially with respect to the cylinder by position engagement with the cylinder. The closure bounds the work space and is sealed with respect to the inner wall by a circumferential seal, but is axially movable with respect to the cylinder when a temperature threshold is exceeded.
2. Description of the Related Art
It is a characteristic of a mono-tube vibration damper that it may be subject to extreme temperature loading, for example, in case of an automobile accident in which the mono-tube vibration damper of the motor vehicle is exposed to fire. In such a situation, the work medium found in the piston device, typically a damping fluid, undergoes a very great expansion, so that a pushing-out force acting on the piston rod and a closure at an axial end of the piston rod overloads the fastening of the closure within the cylinder, the fastening being designed for normal operation. In a worst case scenario, the piston, together with the piston, may be explosively pushed out of the cylinder, due to the pressure level within the mono-tube vibration damper. This can lead to devastating consequences, including injuries to passengers, for example, when the mono-tube damping device is used as a shock absorber in a passenger automobile.
Solutions have been suggested to ameliorate this scenario. These solutions have in common that they attempt to ensure a controlled drop in pressure, rather than an explosion. Some known solutions provide for a deliberately weakened area in the seal to allow for escape of pressurized fluid under extreme pressure situations.
In DE3901449, a solution is proposed in which the seal of the tube is formed so as to have an element that melts upon occurrence of high temperatures, to relieve pressure in the interior. This solution works by providing a deliberately meltable portion of the seal.
U.S. Pat. No. 5,454,455 discloses that at least one bead provides a stop surface for the piston rod-piston unit to prevent its expulsion. However, this solution requires a specially adapted piston and, in some cases, may lead to a loss of useful stroke because the bead is formed within the travel path of the piston.
U.S. Pat. No. 6,840,153, commonly assigned with the present application, discloses a piston-cylinder unit in which a sealing closure is fixed axially by a positive engagement with the cylinder up to a temperature threshold. The closure can carry out an axial movement for relieving pressure above the temperature threshold, in which the relief of the pressure is achieved by providing a cutting device as a tongue-like projection from the wall of the cylinder. Upon axial movement, the cutting device performs a cutting action that cancels the sealing function of the seal. However, this solution requires providing additional structure, i.e., the cutting device, whose only purpose is providing the channel in the high temperature situation.
The danger of uncontrolled explosive axial movement of the piston seal has also complicated and increased the costs of transporting mono-tube vibration dampers. In particular, in some jurisdictions, particularly in the United States, mono-tube vibration dampers are treated, for purposes of transportation, in the same manner as explosives, limiting access to certain roads, bridges and tunnels, increasing the cost of transporting these products. For example, the characteristics discussed above have meant that only certified transportation companies are permitted to transport mono-tube dampers, and only designated roads and border crossings must be used.


Polymerized oils (heat-bodied oils) were hydrogenated to make a semi-crystalline wax-like material. The process of the present invention yields hydrogenated polymerized oils that can be used as bio-renewable replacements for micro-crystalline wax and petrolatums raw materials in products that are comprised of such materials. The present invention is directed to a composition comprising a hydrogenated polymerized oil, a composition consisting essentially of a hydrogenated polymerized oil and a process of preparing a hydrogenated polymerized oil. The present invention is also directed to a composition comprising a hydrogenated blown oil, a composition consisting essentially of a hydrogenated blown oil and a process for preparing a hydrogenated blown oil. The present invention is also directed to a composition comprising a hydrogenated copolymer oil or a hydrogenated copolymer/vegetable oil blend, a composition consisting essentially of a hydrogenated copolymer oil or a hydrogenated copolymer/vegetable oil blend, and a process for preparing a hydrogenated copolymer oil or hydrogenated copolymer/vegetable oil blend.

1. Field of the Invention
Polymerized oils ("heat-bodied oils"), blown oils ("heavy bodied blown oils"), vegetable-based copolymer oils, and vegetable oil/vegetable-based copolymer oil blends were hydrogenated to make a semi-crystalline wax-like material. The hydrogenated products can be used as bio-renewable replacements for micro-crystalline wax and petrolatum raw materials in products that are comprised of such materials.
2. Background Art
Petrolatums and micro-crystalline waxes are components in a large number of products ranging from personal consumer goods to industrial lubricants. Petrolatums and micro-crystalline waxes are derived from petroleum. Thus, the supply of such ubiquitous raw materials is non-renewable. Petrolatums are semisolid mixtures of hydrocarbons, and are neutral, unctuous, odorless, and tasteless. They have applications in the formulation of products such as creams, lotions, hair preparations, cosmetics, candles, ointments, lubricants and coatings. Typical petrolatums have cone penetration of above 100 dmm and less than 275 dmm (ASTM D937). The melting range of petrolatums is about 38° C. to about 60° C. Micro-crystalline wax is a solid mixture of linear, branched and cyclic hydrocarbons derived from petroleum. It is obtained from the heavy lubricating oil fraction derived from crude oil, subsequent to the removal of paraffin wax. Its characteristics resemble those of the natural waxes closely, including its high melting point, high viscosity, flexibility at low temperatures, and high cohesion and adhesion. Micro-crystalline waxes are usually higher in molecular weight, viscosity and melting point than paraffin wax. Typically, micro-crystalline wax melting points range from 54° C. to about 102° C. They have needle penetration of above 3 dmm and less than 100 dmm (ASTM D1321). Viscosities are higher than 5.75 centistokes at 100° C.
Heat polymerized oils ("polymerized oils"), often referred to as heat bodied oils, are prepared from unsaturated oils. Linseed, safflower and soybean oils are commonly used as the starting materials for this process. In addition, fish oils are commonly heat polymerized. Depending on the oil used, the temperature is held between about 288° C. to about 316° C. until a product with a desired viscosity is obtained. Longer reaction times are used to reach a higher viscosity product. The viscosity of polymerized oils can range from P to Z9. During the reaction, the unsaturated triglycerides react to form polymers. As polymerization takes place, new carbon-carbon bonds are formed between triglyceride units.
Polymerized oils have some improved properties for paint, coatings and ink applications in comparison to unsaturated triglycerides. These properties include improved leveling, pigment wetting, and less yellowing. Typical polymerized oils still contain a high amount of unsaturation. The iodine value ("IV") of heat bodied linseed oils ranges from approximately 115-150. Polymerized oils are reactive, viscous liquids at room temperature.
Blown oils differ from polymerized oils. Blown linseed oil is prepared by bubbling air through the oil while heating to temperatures of about 110° C. During the process, the oil is polymerized and partially oxidized.
Vegetable-based copolymer oils such as maleinized and dicyclopentadiene oils are characterized by a fast drying time and water resistance. Blending such copolymer oils with vegetable oils yields oil blends that also possess characteristic properties and provides more diversity of chemical properties. Other attributed properties can include the unique hardness of a dried coating when such copolymer oils are incorporated into paint or coating formulations.

Hydrogenated and partially hydrogenated heat-bodied oils


An integrated circuit device such as an SRAM operating in a battery backup mode, or operating in a quiescent mode when deselected in the operation of a portable electronic device, includes a power dissipation control circuit that reduces the voltage on an internal power supply node so that the memory array is powered at a minimum level sufficient to retain the data stored therein intact.

The present invention relates to integrated circuit devices and more particularly to battery-powered integrated circuits with low power consumption.
Many integrated circuits, such as static random access memory ("SRAM") devices, rely on batteries as backup power supplies to retain the data stored within them while the equipment in which they are used is turned off or the power supply to the equipment has failed. Such integrated circuits are becoming more complex, thus placing greater demands on such backup batteries. Although battery technology is improving, increases in battery capabilities are not keeping up with the increasing power requirements of the integrated circuit devices they serve. There is, accordingly, a need to reduce the power consumption of integrated circuits when used in a battery backup mode.
Additionally, portable electronic devices (computers, cellular phones, etc.) rely upon batteries to supply operational power during general use. Such devices employ integrated circuit devices, which contribute to battery power consumption and reduced operating time between battery rechargings. There is, accordingly, a need to reduce the power consumption of integrated circuits used in portable electronic devices in order to extend the time of operation on a single charge of the battery.

Integrated circuit with power dissipation control


A process is provided for controlling a robotal device, such as a multiaxial industrial robot, by a control unit, with a control core for executing control processes for the robotal device. An interface function checks whether models and/or procedures optionally contained in the control core or additional models and/or transformation procedures and/or special algorithms of kinematic structures, which can be preset at the interface are used as model modules for motion-relevant variables of the robotal device. Special and third-party kinematics can thus also be operated with a control device suitable for executing the process without the control itself having to be modified.

Programmable, computer-aided, i.e., program-based control devices are nowadays usually used for controlling automated robotal devices, such as multiaxial industrial robots, portal robots (robots arranged and movable in a portal frame), SCARA robots (Selective Compliance Assembly Robot Arm—a robot selectively compensating in one plane, pivoted arm robot) or palletizing robots. Procedures for transforming kinematic and dynamic variables and system states of the robot for the purpose of use for different control tasks are a central part of such robot controls. This contains, for example, the position, orientation, velocity and acceleration (of the end effector or the axes), forces and torques in the transmission bearings, desired currents, the kinetic energy of the system, as well as mechanical stresses on the basis of known machine data. These data comprise, for example, geometry descriptions of the robotal device, such as lengths or angles, weights, centers of gravity of masses, and inertias as well as material parameters, such as rigidities, coefficients of friction and coefficients of expansion or the like. Other influential variables may be measured values supplied by sensors, such as temperature, dilatations and distortions.
The transformation of the variables and system states mentioned above as examples may take place in different directions, and a distinction is to be made additionally depending on whether actual or desired values are used as input variables. The following shall be mentioned here as examples:                The determination of the desired current for a robot axis from the desired position of that axis, taking into consideration the effect of the force of gravity, the desired velocity, the desired acceleration and the mass-moment of inertia, taking a desired temperature into consideration;        the determination of an estimated transmission torque from the actual velocity, the actual acceleration and the mass moment of inertia of the axis, taking an actual temperature into consideration; and        the determination of the Cartesian path deviation from desired positions along axes and contouring errors, i.e., differences between the desired position and the actual position.
According to their basic design, the kinematic structures of robotal devices, the so-called "kinematics," are divided into different classes, the division being determined essentially by the number of axes (linear axes, axes of rotation, ball and socket joints, . . . ), the connection of the axes by the components of the robot as well as the relative position of the axes in relation to one another. Whether closed structures ("loops") are formed or whether there is an open kinematic chain is a decisive aspect concerning the connection of the axes. The relative position of the axis is characterized essentially by orthogonality relationships and the intersections of the axes.
The term "class of kinematics" used means that only the above-mentioned embodiment and arrangement of the axes is decisive for the mathematical-physical treatment for the mathematical procedures mentioned, which are to take place in a corresponding control, whereas concrete geometric dimensions of, for example, robots of different sizes for different tasks emerge only as parameters in the corresponding description, without affecting the complexity of the processes per se.
Robot controls are currently typically able to operate robotal devices from only one or a few (fewer than 5) of such classes of kinematics.
The procedures necessary for the control of a robotal device, which were mentioned in the introduction, comprise essentially:
1. Transformations between axis angles and cartesian positions and orientations as well as between the derivations of these variables, i.e., conversion of axis velocities and accelerations into Cartesian velocities and velocities of the reorientation with the corresponding accelerations and vice versa.
These so-called forward and backward transformations are the basis for Cartesian processes; no cartesian path programming of the robot is possible without such procedures. A kinematic model of the robot, i.e., a geometric description without consideration of masses, inertias as well as the resulting forces and torques, is usually used here.
2. Determination of forces, torques and kinetic energy; no time-optimal motion programming and no related model based control is possible without such procedures.
The basis for the calculation is usually a rigid body model (see, among other things, Craig, Introduction to Robotics: Mechanics and Control, Addison Wesley, 1986), wherein friction can also be taken into account. In an expansion, elasticities can also be taken into account in the drivelines. Such a robot model is also called a dynamic model and is the basis for motion planning.
3. Calculation of functions of the kinematic and dynamic variables as auxiliary variables, e.g., of the inertia around a robot axis as an important variable for an axis control.
Variables such as the inertia vary in a robot as a function of other parameters, for example, the payload and axis angles of axes following in an kinematically open chain. Such function values are used, for example, for model-based control.
4. Calculation of partial derivations of the above-mentioned functions.
The partial derivations are needed                for the conversion of variables given for given axes and in a cartesian path by means of the Jakobi matrix;        for the identification of system parameters, e.g., according to the least-squares principle. Such identification methods use partial derivations in the form of the Jakobi matrix. A sensitivity analysis of the parameters requires the Hesse matrix, i.e., the second-order partial derivations. However, such methods are usually used outside the robot control (cf. Kovacs, Computer-aided symbolic robot kinematics, Technische Universität Berlin, 1991).        As additional variables in the control loop.
5. Combinations of the procedures 1 through 3, e.g., kinematic calculation (Procedure 1), taking the deflection of the structure caused by the masses and the force of gravity into account (Procedure 2).
A robot control processes motion commands from different sources, such as program motions, manual displacing motions, and interrupt motions. Motion planning designated the part of a compiler or interpreter that calculates the geometric path of the tool on the basis of certain criteria from the motion commands of a robot programming language (which usually contains the type of point-to-point motion, linear motion, circular motion, spline motion, which are in turn characterized by the target of the motion), possibly including the course of the velocity on the path. This "plan" generated in advance for the course of the motion, is then "interpolated," i.e., points of support are calculated on the planned path—mostly at fixed time intervals—and are passed on—as axis angles—to the drives via the drive interface. For example, a checking is performed in the planning to determine whether the axis angles are within the permissible range (calling of the backward transformation) and/or a time-optimal path is planned (calling of the inverse dynamic model).
The corresponding axis angles are calculated in the interpolation for each cartesian point of support (calling of the backward transformation).
Depending on the current position, velocity, acceleration and load, pre-control torques are calculated in the drive interface. Different derivations are possibly calculated in the control unit.
While no model of the line is used in the classical cascade control of drives and partial derivations are therefore unnecessary, a nonlinear state space model d/dt x=f (x, u) is transformed into a linear model in modern nonlinear, model-based control processes, such as "feedback linearization" (A. Isidori, Non-linear Control Systems, Springer Verlag, 1989) or "backstepping" (Kokotovic: Non-linear and Adaptive Control Design, Wiley Interscience, 1995), and standard methods of linear control theory are then used for this linear model. The transformation contains derivations. This approach of the "Non-linear Model Predictive Control" solves an optimization problem over a certain time horizon in each control cycle (Jan Maciejowski: Predictive Control with Constraints, Prentice Hall, 2001). A zero position of the gradient, which is likewise characterized by partial derivations, is usually determined for the solution of minimization problems.
The above procedures can be performed mathematically with the use of certain standard algorithms of an analytical and/or numeric nature. The following designations will be used below for the mathematical formulation of the procedures:                n number of axes or joints        θi, i=1, . . . , n joint coordinates (angle in case of hinges, lengths in case of rectilinear sliding pairs; scalar)        (x, y, z, a, b, c, S) cartesian position (S: configuration information; see below).
The above-mentioned (transformation) procedures 1 through 5 can be made available in different ways due to the internal structure of the particular robot control used, because they are not set unambiguously in terms of their input and output variables. This shall be shown below based on the example of the backward calculation in case of a kinematics with open chain, i.e., the determination of joint coordinates in case of a given cartesian position of the robot tool (end effector):                Backward calculation in case of a six-armed collapsible robot, i.e., determination of axis angles for a given cartesian position: (x, y, z, a, b, c, S)→(θ1, . . . , θm).
The input variables x, y, z, a, b, c describe the desired position and desired orientation here, for example, via roll-pitch-yaw angles, whereas the additional variable S (typically an integer interpreted as a bit sequence) contains a piece of information that is to be selected from among the different solutions of the calculation. It is implicitly assumed here that an unambiguous solution exists. The additional pieces of information can be preset, for example, in the form of permissible angle ranges.
Backward calculation with quantitative result:(x, y, z, a, b, c)→{N,θi,j}, j=1, . . . , N.
The number of possible solutions N is provided, together with the particular solutions as a result, by this type of backward calculation. The additional information S is not necessary here. It is implicitly assumed that a finite number of solutions exists.                Backward calculation with a result as close to a given position as possible: (x, y, z, a, b, c, φi)→(θ1, . . . , θn).
The solution whose euclidean distance becomes minimal, i.e., ∥φi−θi∥=min, is determined for a given cartesian position (x, y, z, a, b, c) and given axis angles (φi).
In case of processes and devices of the type described in the introduction, which are known in advance, one or more of these different variants are selected and implemented by the manufacturer of the control as a preferred procedure (as preferred procedures). Kinematic data (lengths, angles), dynamic data (masses, inertias) and optionally other physical parameters are stored as parameters for the procedures 1 through 5 in corresponding machine descriptions that belong to them.
A fundamental restriction arises during the implementation for mathematical reasons: It is not possible to specify a general mathematical process for determining all procedures 1 through 5 from a given geometric and dynamic description of a robot. Some fundamental mathematical problems in this connection are:                The backward calculation in case of kinematics with open chain can be reduced to the determination of the zero positions of polynomials via the so-called tan/2 substitution, but no analytical solution exists for general polynomials beginning from the fifth degree (van der Waerden, Algebra I/II, Springer, 1993). Thus, no analytical solution of the backward calculation is possible for any desired oblique positions of axes and for any desired distance between axes. In addition, the effort needed for the calculation also increases greatly for the analytical solution—if any such solution exists—in case of the use of the most modern methods (Kovacs, elsewhere). An analytical solution is to be preferred to a numerical one, because numerical methods are, for example, incapable of determining the number of possible solutions.        The forward calculation is trivial in case of kinematics with open chain, because only matrices must be multiplied here (Craig, elsewhere; Paul, Robot Manipulators, MIT Press, 1983). No analytical forward calculation is known, as a rule, in case of kinematics with kinematic loops, e.g., hexapods, so that what was said above in connection with the backward calculation correspondingly applies here. The backward calculation in case of loop kinematics can, by contrast, be often determined in a simple manner analytically (Möller, A Method for the Automatic Analysis of Multi-Loop Spatial Mechanisms, Dissertation, Universität Stuttgart, 1992).        The calculation of the "inverse dynamics" is possible in case of open kinematic chains via the so-called Euler-Newton algorithm, if the axis coordinates are known. As in the case of the backward calculation of kinematics with open chain, there is no general solution in case of closed kinematic chains.        If no analytical calculation method is known in case of a certain procedure, there is usually no possibility for analytical calculation for the partial derivatives, either (see above).
Nevertheless, the manufacturer of a control strives to be able to operate the largest possible number of types of kinematics with a given control. This is especially true if this manufacturer offers both controls and kinematics. On the other hand, a control shall also be able to be operated with kinematics of third-party manufacturers, who do not offer controls of their own themselves. A manufacturer of controls had two possibilities for this up to now: On the one hand, the implementation of all procedures 1 through 5 for a limited number of kinematics, and, on the other hand, the use of (iterative) numeric methods combined with analytical solutions. The fact that all procedures 1 through 5 must be implemented for each kinematic, so that the scope of the control also grows with the number of the kinematics to be supported at will, besides the effort needed for the development and care, is to be considered to be disadvantageous here, in particular. Consequently, if newly developed kinematics or third-party kinematics of another manufacturer are to be operated with a given control, this is possible only if the corresponding kinematics falls within one of the already implemented classes of kinematics or can be calculated with the procedures of these kinematics.
The use of (iterative) numerical methods combined with analytical solutions shall be explained below based on the interaction of forward and backward transformations in case of a kinematic with open chain. Only a matrix multiplication is necessary for the forward calculation, and the matrix entries are composed of the joint variables and kinematic descriptions (for example, the Denavit-Hartenberg parameters) as well as of the coordinate systems for the tool, the robot base, etc. The position of the end effector and the orientation P of the end effector are given mathematically as a function of the axis variable θi: P=f(θi).
The necessary parameters can be stored in a data set for the machine description. According to the statements made above, the back transformation cannot be solved, in general, analytically, so that a method for the numerical solution of systems of equations of the above form is used, instead, for a given P (Münch, Universal coordinate transformation for industrial robots, www.maschinenbau.hs-magdeburg.de/personal/bargfrede/-fue_robotics.html).
What is to be considered to be disadvantageous in this procedure is, in particular, the fact that it is not possible to implement the procedures 1 through 5 in a generally valid form, and especially the lack of ambiguity of the solution, as well as the convergence characteristic, which is critical for real-time applications, and the often insufficient accuracy argue against the numerical methods, which are used by necessity.
A manufacturer of kinematics nowadays depends on software development cycles of the control manufacturer: Even if the manufacturer of the kinematics is able to implement the procedures 1 through 5 by calculation, he has no possibility of utilizing these with an existing control software, Moreover, the manufacturer of the kinematics usually wishes to develop the procedures 1 through 5 himself in order not to be forced to have to disclose operational know-how to the manufacturer of the control. However, this is not possible in case of the processes and the devices of the type mentioned in the introduction, because the manufacturer of the control would conversely have to disclose detail information concerning the control to the manufacturer of the kinematics.

Process and device for controlling the robotal device


Disclosed are automatic vertical moving systems and control methods therefor for regularly lowering an object with predetermined depth and time intervals to a desired maximum depth in water and then raising it to the surface of water to analyze vertical changes of water quality in lakes and reservoirs continuously, comprising a barge for floating on the surface of the water so as to be retained at a desired position thereon, a winch mounted on the barge, rotatable in opposite directions to lower and raise the object and provided with a wire rope wound and rewound around the winch according to the rotational directions, the wire rope being connected at its free end to the object, an electric motor connected to the winch for rotating in opposite directions, and a motor controller for controlling the motor to be regularly activated and deactivated and to be changed in rotational directions thereof. Also, an assistant line may be further included to prevent entanglement of the wire rope and possible loss of the object due to external impact of unusual matters in water.

1. Field of the Invention
The present invention relates generally to automatic vertical moving systems and control methods therefor, and more particularly to automatic vertical moving systems and control methods therefor for regularly lowering an object with predetermined depth and time intervals to a desired maximum depth in water and then raising it to the surface of water.
2. General Background
Recently, automatic water quality analysis system with associated internal memory, such as YSI6000 from YSI Inc., U.S.A. became commercially available to monitor water quality in lakes, reservoirs, rivers and coastal areas. The analysis system is generally provided with sensors for automatically measuring water temperature, dissolved oxygen concentration, conductivity, pH, salinity, turbidity, oxidation reduction potential, and etc., and built-in memory to store data collected during operation.
The analysis system, however, has no moving capability. So, for observing the vertical changes of the water quality, especially in lakes and reservoirs, the system should be lowered and raised manually, which takes a long time and is inefficient to acquire the desired data.

Automatic vertical moving systems and control methods therefor

508-304 Expo Apt., Jeonmin-dong, Yusung-gu


A rotary machine comprises a hub including: a blade-facing hub portion including a first blade-facing surface which faces a spherical hub-side end surface of a variable blade and which has a first spherical region having a spherical shape; an upstream hub portion disposed upstream of the blade-facing hub portion in an axial direction of the hub and having a first outer-peripheral surface being adjacent to the first blade-facing surface in the axial direction; and a downstream hub portion disposed downstream of the blade-facing hub portion in the axial direction and having a second outer peripheral surface being adjacent to the first blade-facing surface in the axial direction. At least one of following condition (a) or (b) is satisfied: (a) a downstream end of the first outer peripheral surface is disposed on an outer side of an upstream end of the first blade-facing surface in the radial direction of the hub; (b) an upstream end of the second outer peripheral surface is disposed on an outer side of a downstream end of the first blade-facing surface in the radial direction of the hub.

In a rotary machine such as a compressor and a turbine, at least one of a stationary vane or a rotor blade may be configured as a variable blade that is revolvable about a pivot axis along the radial direction of a hub, to adjust the attack angle with respect to flow.
In a rotary machine provided with such a variable blade, if the variable blade is configured such that the hub-side end surface of the variable blade does not interfere with the blade-facing surface of the hub in the rotation range of the variable blade, clearance between the hub-side end surface of the variable blade and the blade-facing surface of the hub is likely to increase when the variable blade is revolved toward the close side (in a direction the angle between the chord line of the variable blade and the axial direction of the hub increases). If the clearance between the hub-side end surface of the variable blade and the blade-facing surface of the hub increases, loss due to a leaking flow that passes through the clearance (hereinafter, described as clearance loss) increases, and the efficiency of the rotary machine may decrease.
Patent Document 1 discloses a rotary machine with a variable blade having a spherically-shaped hub-side end surface recessed outward in the radial direction of the hub and a blade-facing surface of a hub having a spherically-shaped spherical region protruding outward in the radial direction of the hub, so that the clearance does not increase at rotation of the variable blade toward the close side.


Disclosed is a lighting optical system, including: a lighting optical system, comprising, a light separating means for separating incident light into three colored lights according to wavelength bands, and progressing the three separate colored lights for the same amount of distance from a light source of the incident light through different optical paths, and emitting the three separate colored lights to a first, a second and a third LCD, a first, a second and a third LCD for respectively receiving the three colored lights from the light separating means, and forming images therefrom; and a light combining means for synthesizing the three colored lights emitted from the first, second and third LCD, and outputting the synthesized light outside.

1. Field of the Invention
The present invention relates to a lighting optical system, and more particularly, to a small-sized lighting optical system with an improved performance.
2. Background of the Related Art
Technical advances in display devices have created lighter, thinner and larger screens.
Examples of such display device include projectors or big-screen projection TVs.
In recent years, micro devices have been introduced as a new technology for projectors and projection TVs. Examples of micro device include backlit LCD (HTPS) panels, reflective LCoS (Liquid Crystal on Silicon) panels, and DMD (Digital Micromirror Device) panels.
Depending on how many micro devices are used, the projectors or the projection systems can be categorized as single-panel, double-panel, and triple-panel projectors or projection systems.
In case of the backlit LCD (HTPS), triple-panel optical systems are widely used. As for the LCoS triple-panel, double-panel and single-panel optical systems are all available. As for the DMD, single-panel optical systems are generally used. However, it should be noticed that each of these examples does not always have advantages only.
For instance, a triple-panel optical system for the backlit LCD (HTPS) has too many optical elements and includes a relay lens for compensating optical path differences.
Besides the relay lens, the backlit LCD uses an X-prism for combining colors.
Since there is an optical path difference among R, G and B lights, there was a need to develop a new optical system capable of compensating the optical path difference. That was how the X-prism got involved therein.
FIG. 1 is a schematic diagram illustrating a triple-panel projection system according to a related art.
Referring to FIG. 1, the triple-panel projection system includes a lamp 110 for emitting light; fly-eye lenses 120, 121 for splitting the light emitted from the lamp 110 into micro lens cell units; a PBS (Polarizing Beam Array) 130 for forming incident light into a linearly polarized light with one optical axis; condensing lenses 140, 142, 142 for condensing light; dichroic lenses 151, 152 for splitting the light into R, G and B colors; LCDs (Liquid Crystal Displays) 161, 162, 163 for providing or displaying R, G and B images; mirrors 171, 172, 173 for changing a traveling path of each of the split lights by the dichroic lenses 151, 152 so that the split lights travel to their corresponding LCDs 161, 162, 163, respectively; relay lenses 181, 182 disposed on the traveling path of light to focus the light to a desired position; and an X-prism 190 for combining images from the LCDs 161, 162, 163, respectively.
The following will now explain the operation of the related art projection system described above.
The light from the lamp 110 incidents on the first dichroic lens 151.
The first dichroic lens 151 then reflects red light and transmits cyan light.
The red light reflected from the first dichroic lens 151 is reflected again by the mirror 171 and incidents on the first LCD 161.
The second dichroic lens 152 reflects green light and transmits blue light.
The green light reflected from the second dichroic lens 152 incidents on the second LCD 162.
The blue light transmitted through the second dichroic lens 152 goes through the relay system 200 consisting of the first and second relay lenses 181, 182, and the first and second mirrors 172, 173, and as a result the light path difference of the blue light from the red and green lights is compensated. Afterwards the blue light incidents on the third LCD 163.
The incident lights on the first, second and third LCDs 161, 162, 163 experience phase modulation according to an input signal and put corresponding image information therein. The lights containing the image information are then emitted from the first, second and third LCDs 161, 162, 163, are synthesized through the X-prism 190, and reach a screen (not shown) by means of a projection lens (not shown).
The above-described projection system (or projector) requires a number of optical elements for splitting and combining colors.
It also requires the relay system for compensating the light path of the blue light to make it equal to those of the red and green lights.
Because of the relay system the size of a light engine is increased, more optical elements are required, and a great number of elements should be adjusted to array an optical axis.
Meanwhile, the red and blue lights were in the same polarization state to obtain a maximum light efficiency taking advantage of the X-prism, but the green light had a different polarization state. Therefore, a retardation plate was additionally installed next to the X-prism to make the green light have the same polarization state with that of the red and blue lights.


A mobile device to supply information, notably for tourism, to a user, has a storage device that is fragmented into sectors, each sector corresponding at least to one given geographic location and showing an address, a tracker of said device, capable of delivering positioning information, corresponding to the location wherein the device is situated, a processor capable of drawing up a reading instruction containing at least the address of the sector of the information to be supplied in relation the positioning information on the basis of one or several correspondence tables.

When traveling, when one wishes to obtain information regarding the location, one may first of all look up literary tourist guides. Nevertheless, this solution is rather impractical since it implies turning the pages of the guide constantly and dedicating some time to reading said guide.
To remedy these shortcomings, It has been suggested to use audio or video tape players in which one has loaded tapes with pre-recorded comments on the locations encountered during said travel.
It is obvious that such a solution may be cumbersome for the user who would then have to stop and switch his player on again constantly to avoid any delay between the information provided and his progress on the road. Consequently, as the driver has heard or seen comments too early or too late with respect to the speed of his progress, he does not know where he is any longer and gives up using said device.
It should be noted that, with the solutions mentioned above, the nature of the information available remains limited to the language and/or the subject chosen (history, architecture, gastronomy, leisure, . . . ).
Besides, in the guiding technology, vehicle guiding devices are already known, which enable positioning of said vehicles and/or recommending a particular route to the driver.
Nevertheless, such systems are developed for vehicles only and solely regard the route followed by said vehicle.

Portable device supplying tourist information


A tuning method and a tuning apparatus for tuning a filter are disclosed. The tuning method includes: configuring the filter as a VCO; utilizing the VCO to generate an oscillation signal according to a driving signal; comparing a frequency of the oscillation signal with a reference frequency to generate a comparison result; converting the comparison result into the driving signal in order to establish a feedback mechanism. Therefore, the inner components such as the gm and capacitance inside the VCO are completely tuned when the VCO generates an oscillation signal having a wanted frequency. Since the VCO is inside the filter and the components of the filter and the VCO are similar, the driving signal can be utilized to make the filter operate in a desired center frequency under a well-designed relationship between the frequency of the oscillation signal and the center frequency.

The invention relates to a tuning apparatus and method, and more particularly, to an apparatus for tuning a center frequency of a filter and related method thereof.
Filters have been widely used in many applications. In general, the implementations of the filters can be roughly divided into two types of structures, which are discrete-time switch-capacitor filters and continuous-time filters (including gm-C filter, MOSFET-C filter, etc.). Because of the limitation of clock frequency, a high-frequency filter is often implemented by the continuous-time structure. Furthermore, in different types of continuous-time filters, the gm-C filter is the most general.
However, the variation degree of some characteristics, such as cut-off frequency (center frequency), of the continuous-time filter is often larger than 30% due to the influences of variations in the manufacturing process and temperature variations. Therefore, a tuning mechanism should be added into the filter to overcome the process and temperature variations such that the frequency response of the filter is not affected by the influences of the process and temperature variations.
In the gm-C filter structure, the cut-off frequency (center frequency) can be represented by the following equation wc=K*gm,u/C, where gm,u represents a gm(trans-conductance) value of a trans-conductor per unit, C represents a total capacitance value corresponding to a node, and K is a scaling factor larger than 0. From the above equation, it can be seen that when the center frequency wc deviates from a target value, the gm value gm,u or the capacitance value C can be adjusted to tune the center frequency wc back to the target value. In order to achieve the purpose of tuning the center frequency wc, either the gm value gm,u or the capacitance value C can be adjusted. Please note that adjusting the gm value gm,u or the capacitance value C is equivalent. The spirit of the two adjusting methods is the same.
Taking the method of adjusting the gm value gm,u as an example, please refer to FIG. 1, which shows a conventional tuning structure 100 of adjusting the gm value gm,u of a main filter 110. Please note, the main filter 110 is a target filter to be tuned. Furthermore, in FIG. 1, the tuning operation is performed by a PLL (including FD 120 the charge pump 121 and the loop filter 122) cooperating with a VCO 130. Please note, under the tuning structure 110 shown in FIG. 1, the VCO 130 is better composed of the same trans-conductor cells as those of the main filter 110, the VCO 130 has the same environment (e.g., loading, etc.) as the main filter 110, and the gm value of the trans-conductor circuits of the VCO 130 and the main filter 110 are controlled by the same control signal Vc. Therefore, if the tracking relationship between the VCO 130 and the main filter 110 is better, the tuning operation of the tuning structure can be more accurate. In other words, when the VCO 130 is tuned, the main filter 110 is also tuned because they have similar environment. Assume that the center frequency wc of the main filter 110 is ideally equal to the K*gm,u/C, and the oscillation frequency wo of the VCO 130 is equal to N*gm,u/C. When the PD 120 of the PLL is locked to a certain frequency, the control signal Vc is adjusted to change the gm value gm,u such that the fo=(½π) (N*gm,u/C)=fref. As mentioned previously, the main filter 110 and the VCO 130 has a good tracking relationship (e.g., they have the same gm value gm,u). Therefore, fc=(½π) (K*gm,u/C)=(K/N) fo=(K/N) fref. Obviously, if the values K, N, and fref can be selected properly, the center frequency of the main filter 110 can be tuned to a target frequency.
Please refer to FIG. 2, which is a diagram of another conventional tuning structure 200. Please note, the tuning structure 200 shown in FIG. 2 utilizes a similar concept. The tuning structure 200 utilizes similar trans-conductor cells to form the master filter 230 (in general, the master filter shown in FIG. 2 has lower levels), and utilizes the characteristic of the master filter 230 to perform the tuning tasks. For example, a two-level Biquad LPF has a 90 degree phase delay at the point where wo=N*gm,u/C. Therefore, when the signal having the frequency fref is inputted into the master filter 230, the entire tuning structure 200 utilizes the phase detector (PD) 220 to determine whether the phase difference is 90 degrees. Additionally, recall as mentioned previously, the feedback mechanism is implemented by a PLL including a charge pump 221 and a loop filter 222, the negative feedback mechanism adjusts the control signal Vc to make the fo=(½π) (N*gm,u/C)=fref. From the above-mentioned structure, it can be easily seen that fc=(½π) (K*gm,u/C)=(K/N) fo=(K/N) fref. Therefore, the tuning structure 200 shown in FIG. 2 can also achieve the same tuning goal.
The above-mentioned structures both needs a PLL including a PD, a charge pump, and a loop filter. It is known that the PLL occupies a larger area and as one result, this increases the cost. Please refer to FIG. 3, which is a diagram of another conventional tuning mechanism 300. The tuning mechanism 300 utilizes a digital circuit 320 to perform a negative feedback control. The entire tuning method shown in FIG. 3 is more similar to the tuning structure shown in FIG. 1. The difference between the tuning structures shown in FIG. 3 and FIG. 1 is that the tuning structure shown in FIG. 3 utilizes a digital circuit 320 (i.e., a digital FD) to compare the frequency fref with the frequency fc generated by the VCO 330 instead of utilizing a PLL. Thereby, the comparison result is transformed into a control signal through a DAC 340 in order to adjust the frequency fc. Similarly, because of the tracking relationship between the VCO 330 and the main filter, when the oscillation signal of the VCO 330 is tuned, the cut-off frequency of the main filter 310 can be tuned successfully.
The influences caused by the process and temperature variations upon the frequency fc, can be alleviated through the above-mentioned tuning mechanisms. Obviously, the above-mentioned tuning mechanisms need either a VCO or a master filter. Furthermore, either the VCO or the master filter is often a two-level system. In addition, in order to make the environment similar to the main filter, all dummy devices, dummy loading, and other circuits, which have originally been set up in the main filter, also have to be copied and implemented inside the tuning structure (e.g., the above-mentioned VCO or master filter) to make the environment similar. In most of the applications, the tuning structure often occupies a huge area larger than 20% of the entire circuit. Therefore, the above-mentioned tuning mechanism consumes a large area and high cost resulting in an uneconomical solution.

Apparatus and method for tuning center frequency of a filter


A sleeve system for lining a hole of a surgical screw in a bone or tissue of a subject includes a woven sleeve having a first end, a second end, and an elongated body. The elongated body has an interior portion. The first end includes an aperture that accesses the interior portion. The sleeve system also includes a protective tip that may be inserted through the aperture of the first end of the sleeve. The protective tip at least partially protects tissue of the subject from a screw positioned within the sleeve.

U.S. Pat. No. 7,758,642, which is incorporated herein in its entirety by reference, discloses the use of a perforated sheet to assist in implant fixation. U.S. Pat. No. 5,984,926, issued on Nov. 16, 1999, which is incorporated herein in its entirety by reference, discloses a suture based device for saving screw holes in bone, as well as various problems associated with achieving adequate grip or achieving adequate bone screw purchase strength in cancellous bone or osteopenic/osteoporotic bone or bone where the fractures or poor quality bone preclude optimal engagement of bone.
Because of the variety of types of injuries, bone conditions and anatomical variations and limitations, there remains a need for devices that improve screw purchase strength in a variety of clinical situations. Known methods, such as the method described in U.S. Pat. No. 8,163,031 to Truckai et al., incorporated herein in its entirety by reference, discloses a sleeve for improving the grip of the screw with the bone. However, the Truckai et al. device does not allow the surgeon to adjust the thickness of the insert. Also, a significant problem associated with pedicle screws is that they often move within the screw hole, and the sharp point of the screw causes damage to the bone around the tip, as well as causes pain to the patient due to the close proximity of the screw tip and nerves or other sensitive tissue.
Still further, the device disclosed in U.S. Pat. No. 5,984,926 may be inserted using a rod, and thus the device would benefit from the use of an insert that would assist the surgeon in engaging and centering the sleeve while pushing the sleeve into the screw hole.

Surgical screw hole liner devices and related methods


A damping valve includes: a ring-shaped leaf valve having either one of an outer circumference and an inner circumference as a free end, the free end being allowed to be deflected toward both sides in the axial direction; a ring-shaped opposing portion that opposes the free end of the leaf valve with a gap; a first sub leaf valve stacked on one side of the leaf valve in the axial direction; and a first passage formed in the leaf valve so as to extend in parallel with the gap, the first passage being configured to be opened when the leaf valve is deflected in the direction away from the first sub leaf valve.

Conventionally, valves are used, for example, for generating a damping force by imparting resistance to a flow of liquid that is caused during extension/contraction of a shock absorber. As such valves, there is a known valve that is provided with: a leaf valve, in which the one of an inner circumference and an outer circumference is formed as a fixed end and the other thereof is formed as a free end that is allowed to be deflected in both sides in the axial direction; and an opposing portion with a ring-shaped opposing surface that protrudes toward an end surface of the free end of the leaf valve so as to form a minute gap with the end surface of the free end in a state in which the opposing portion opposes to the end surface (for example, see JPH 2-76937A).
In the above-described configuration, because the minute gap is formed between the end surface of the free end of the leaf valve and the opposing surface of the opposing portion, even in a state in which the end surface and the free end are opposed to each other, an end portion of the leaf valve on the free end side can move in both sides in the axial direction from a position at which the end surface thereof opposes to the opposing surface of the opposing portion to a position at which the end surface and the opposing surface are no longer opposed to each other. With such a configuration, as an extension/contraction speed (piston speed) of the shock absorber is increased, the end portion of the leaf valve on the free end side is deflected such that the end surface thereof and the opposing surface of the opposing portion come to no longer be opposed to each other. As a deflected amount of the leaf valve is increased along with the increase in the piston speed, the gap formed between the free end and the opposing portion is increased to increase a flow amount of liquid passing through the leaf valve.


Disclosed is a method for producing biodegradable polymer and ethanol converted from carbon source by using recombinant microorganisms, comprising the steps of: (A) providing recombinant microorganisms transformed with plasmids containing at least a gene encoding for glycerol utilizing enzyme and a gene encoding for polyhydroxyalkanoate synthase; (B) culturing the recombinant microorganisms in a medium containing glycerol; (C) inducing expression of the genes of step (A), thereby obtaining polyhydroxyalkanoate and ethanol; and (D) recovering the polyhydroxyalkanoate and the ethanol; wherein the recombinant microorganisms have a glycerol utilization rate more than 90% (w/w), and have polyhydroxyalkanoate accumulated therein to a biomass content thereof at least 30% (w/w).

1. Field of the Invention
The present invention relates generally to a method for producing particular products by recombinant microorganisms, and more particularly, to provide a method for producing biodegradable polymer and biomass fuel converted from carbon source by using recombinant microorganisms.
2. The Prior Arts
With progressing development of bio-diesel, using B2 bio-diesel blend as fuel of vehicles was promoted in Taiwan from 2000. However, the increasing of discarded crude glycerol accompanied with bio-diesel production results in environmental pollution. Although the discarded crude glycerol could be recycled and refined, the economic benefits are limited. Thus, treatment of this crude glycerol, byproduct of production, has becoming the key point for the development of bio-diesel. Regarding the trends of the world, most advanced countries have related laws made for energy policy, and majorities of which have being as principal items for national development. For example, the vehicles using petroleum must be replaced by those using substitute fuels with designate percentage in a given time; as well as increasing the percentage thereof gradually. Under those trends, the production quantities of bio-diesel will be increased relatively, and at that time it will certainly be a serious problem inasmuch as the enormous discarded crude glycerol. Therefore, if the crude glycerol can be efficiently processed and utilized by microorganisms, for example, producing biodegradable polyhydroxyalkanoate (PHA), the additional value resulted from bio-diesel production will be enhanced consequently. PHA, a kind of polyester existing in microorganisms, is generally served as a carbon and energy source. When the microorganisms are in a nutrients-limit condition, they will generate a hydrophobic inclusion body encompassing the PHA within the cell. The hydrophobic granules may have the PHA accumulation content over 90% of their dry weight under certain condition. Concerning the chemical structure, PHA is a kind of polymer composed of hydroxyl fatty acid, i.e., hydroxyalkanoic acid (HA); the structural formula thereof is presented as follows:
wherein the molecular weight range generally from several ten thousands to several millions, m is 1-4, n is 100-3000, and R is C1-C5. Variable kinds of PHA and monomers thereof differ in R group; in general, PHA compound comprising polyhydroxybutyrate (PHB), polyhydroxyvalerate (PHV), and copolymer of PHB and PHV (i.e., linked both with PHB and PHV), etc. With different kinds of monomer composed, PHA may not only be a kind of hard and brittle plastics but also a soft elastomer, which involve a series of different characteristics thereof. PHA could be synthesized by using bin-renewable resources isolated from organisms as raw materials; moreover, it could be degraded by organisms, e.g., bacteria, when existing in natural environment, and hence could be a plastic substitute for conventional undegradable plastic. As a biocompatible and biodegradable thermoplastic polymer, it is suggested that PHA will have a good application in medical science and industry.
The production and development of PHA were started in the 1970s. ICI Company produced PHA by using microorganisms in the soil through fermentation process at that time. Meanwhile, Massachusetts Institute of Technology (MIT) introduced engineering techniques to produce PHA by using microorganisms. Owing to the research results, Metabolix Company was derived and established. The patents concerning the abovementioned results were assigned to Zeneca Company, and subsequently assigned to Monsanto Company. Metabolix Company, having purchased related patents from Monsanto Company in 2001, researched and developed on the basis of the techniques from patents and their own, and cooperated afterward with Archer Daniels Midland (ADM) to set up bulk scale production of PHA achieving 50 thousand tons in 2004 and intending to put to production in 2007-2008. The cost for PHA production is US$3-5/kg by fermentation method pertain to Metabolix Company; however, the cost may be reduced to less than US$2-3/kg by industrial process. Therefore, Metabolix Company planned to further abate cost to US$2-3/kg by improving innovative devices. Nevertheless, the highest part of cost is about the medium components employed for fermentation, which accounts for 50-80% of the operation cost. By the analysis of upstream and downstream material cost, the cost of carbon source and nitrogen source is more than 80% of total material cost. As the high cost of carbon source, commercialization and competition of PHA products are obstructed. Thus, if the crude glycerol produced by the bio-diesel factory could be recycled and utilized as a carbon source, it will be beneficial to achieve the objective of cost reduction for PHA production.
On the whole, the national research concerning PHA includes the strain improvement (e.g., Escherichia coli, Wauteria eutropha), gene mutation (e.g., genes of phaP, phaR), scale-up of culturing volumes, processes for fermentation production, tests for industrial application etc., what is more, the gene research and the strain improvements are both emphasized. In respect of Mainland China, which adopts biodegradable biopolymer material as the main issue for environmental protection of 2008 Beijing Olympic Games; meanwhile, the achievements regarding biodegradable plastics research are also planned to propagandize intensively. Besides, in Korea, the technique of high cell density fermentation for PHA production does have a referring value; moreover, the aspects of bacteria species refinement and the pathway improvement of metabolism also have been studied deeply. Furthermore, the German research group decodes the genomic DNA of Ralstonia eutropha, where the homologues genes of phaC, phaB, and phaA are found. It is not only a crucial breakthrough research for Ralstonia eutropha, but also could be applied to improve the production of other bacterial species by gene alignment and genetic engineering. In addition, various types of bio-polyester pertained to PHA, e.g., 3-hydroxybutyrate (3HB), 3-hydroxyvalerate (3HV), and 4-hydroxybutyrate (4HB), are produced by using Bacillus cereus SPV from the research group of England.
Furthermore, the research group of Switzerland puts emphasis on using genetically modified plants, being different from microbiological expression system, to produce PHA. In U.S.A., other than the studies of bio-polyester production by means of recombinant E. coli expression system, the characteristics of copolymer are also analyzed so as to amplify the downstream application fields. Additionally, the research group of Australia deeply research on the gene regulation system of PHA production in bacteria pertained to Pseudomonas family, which indicates that the bacteria even of the same family present significant differences in gene level expression. Moreover, the Indian research group uses photosynthesis bacteria for PHA production; and Canada tests the biodegradability of PHA for further research.

Method for producing biodegradable polymer and biomass fuel converted from carbon source by recombinant microorganisms


The invention provides an optical tracing system for measuring elevation of a reference body above a reference surface, comprising a laser transmitter radiating a light beam in a downward direction; an optical unit for imaging the radiation reflected from said body and said surface on a detector; filtering means for separating light originating from said laser from background illumination, and signal processing means for identifying said reference body and for measuring the elevation thereof from said surface.

The present invention relates to an optical tracing system for measuring elevation of a reference body above a reference surface. The system of the invention is especially useful when attached to the blade or another part of an earth-moving, paving or grading machine which is being driven at normal working speed. The system detects a reference body such as a string line, cable, rod, curb or surface, measures and memorizes data regarding the initial elevation of the reference body above the ground, and thereafter continuously measures and compares measured values with the memorized data.
Elevation measurements are commonly performed in road construction operations and an ultrasonic range finder is normally used for this purpose. However, such a device requires some care in order to ensure accurate measurements, as it is adversely affected by wind, air temperature and temperature gradients. In addition, ultrasonic waves are reflected not only from the measured body, but also from the background, such as the ground beneath the body, which results in faulty measurements.
Optical tracers, based on the triangulation measuring technique, have been proposed in the past, e.g., in U.S. Pat. Nos. 4,456,826; 4,774,403; 5,129,725; 5,015,868 and 4,911,551. The first of these patents discloses a device based on rotating multi-facet mirrors, which is expensive and not durable if utilized in an open-air environment. In the other patents, a CCD array serving as a detector and a laser utilized as a light source are described. Such components are hardly compatible with road construction conditions, which, in general, require the low laser power dictated by eye safety regulations. These conditions include a large range of operation temperatures, below -20.degree. C. in winter and up to 55.degree. C. in summer in direct sunshine, and an extremely wide range of outer illumination conditions, from cloudy and rainy weather to bright sunshine, and also require the ability to measure elevation above different reference bodies (string line, cable, rod, curb, ground or passageway) wherein the type of reference body has to be automatically identified, as well as limited size and power consumption of the device used.

Optical tracing system for measuring elevation of a reference body above a reference surface


The liquid chromatograph is provided with: a first dimension analysis unit; a diluent determination unit configured to determine the composition and flow rate of a diluent for diluting an eluent flowing out of the first dimension analysis unit based on the composition of the eluent flowing from the first dimension analysis unit; and a diluent liquid delivery unit configured to deliver the diluent of the composition determined by the diluent determination unit at the flow rate determined by the diluent determination unit so as to be delivered to the trap column while diluting the eluent flowing from the first dimension analysis unit.

In a second dimension high-performance liquid chromatography, an eluent containing one or more analysis target components is fractionated on a downstream side of a column in a first dimension analysis, and the analysis target components are adsorbed on a trap column, and then the analysis target components adsorbed by the trap column are released to be guided to another column to perform a second dimension analysis (see Patent Document 1).
In such a high-performance liquid chromatograph, for the purpose of preventing the precipitation of the separated high concentration analysis target components in the eluent in the flow path, it is configured such that a large amount of a mobile phase containing a constant concentration organic solvent is flowed into the flow path to dilute the eluent. Also, when introducing an eluent to a trap column to adsorb it, when the organic solvent concentration in the eluent is high, it may occur that the analysis target components are not adsorbed by the trap column. For this reason, in order to prevent such a situation, it is performed to lower the organic solvent concentration of the eluent by using a diluent.


Disclosed are anode electrode structures for microbial fuel cell (MFC) devices, systems and methods for treating wastewater and generating electrical energy through a bioelectrochemical waste-to-energy conversion process. In some aspects, an anode electrode includes a conductive core and a plurality of sheets of conductive textile material wound around the conductive core. In some aspects, the anode electrode is produced by cutting sheets of a conductive textile material to form a stem and a plurality of branches connected to the stem. The conductive textile material is pretreated to enhance the surface area, hydrophilicity, microbial attachment, and/or electrochemical activity of the conductive textile material. The sheets are stacked together and wound around a conductive core to produce the anode electrode. In implementations, the anode electrode can be used to transfer electrons removed from wastewater surrounding the branched electrode via an oxidation reaction on the electrode surface within the an MFC device.

A microbial fuel cell (MFC) is a bioelectrochemical system that uses living cells, such as bacteria, and mimics and/or promotes their natural interactions to produce electric current. Some example MFC systems include mediated MFCs, which use a mediator for transferring electrons from the bacteria cell to the anode. Other MFC systems include unmediated MFCs, which utilize types of bacteria that typically have electrochemically active redox proteins (e.g., cytochromes) on their outer membrane that can transfer electrons directly to the anode.

Branched anodes for microbial fuel cells


An optical transmission module of the invention uses a semiconductor chip forming an optical modulator integrated laser on a semi-insulating semiconductor board. An input transfer line and an anode electrode of an optical modulator element are connected by a first bonding wire. The anode electrode of the optical modulator element and one of the ends of a terminal resistor element are connected by a second bonding wire. A cathode electrode of the optical modulator element and the other end of the terminal resistor element are connected by a third bonding wire. The cathode electrode of the optical modulator element and a ground electrode are connected by a fourth bonding wire. A joint portion between the first bonding wire and the input transfer line is arranged on an opposite side to a joint portion between the fourth bonding wire and the ground electrode while interposing the semiconductor chip between them.

This invention relates to an optical transmission module for optical communication, and particularly to an optical transmission module used for a transmission portion of an optical transceiver having a high transfer rate (10 Gbit/s, for example).
An optical transmission module using a semiconductor laser is one of the key devices of transceivers for optical fiber transfer. An operation speed of optical transmission modules has been improved with recent popularization of broadband networks and optical transmission modules having a transfer rate of up to 10 Gbit/s have now been used widely. The optical transmission modules suitable for this application are strongly required to possess small size and low cost and to achieve high transmission waveform quality.
Optical transmission modules directed to improve high frequency characteristics are disclosed in JP-A-2001-257412 (corresponding to U.S. Pat. No. 6,823,145) and JP-A-2001-308130 (corresponding to U.S. Pat. No. 6,507,111). These patent references JP-A-2001-257412 and JP-A-2001-308130 describe the optical transmission modules that simultaneously reduce a small signal reflection coefficient (S11) at a characteristic impedance of 50 Ohms on the high frequency input side and secure a 3 dB band in a small signal passage characteristic (S21) of an optical modulator by optimizing the relation of inductances between a first bonding wire for connecting a modulator and a signal line of an electro-absorption modulator type optical modulator integrated laser diode and a second bonding wire for connecting the modulator and a terminal resistor.
As disclosed in JP-A-2001-257412 and JP-A-2001-308130, one of the ends of the terminal resistor is connected by the bonding wire to the modulator but the other end is connected to a ground electrode through a ground via-hole, or the like, without being connected to the modulator in the optical transmission module according to the related art. A cathode electrode or anode electrode of an electro-absorption type modulator is disposed on the chip back surface of the modulator integrated laser diode.
JP-A-2005-353910 (corresponding to US Patent Publication No. 2005/275053) describes, on the other hand, an electro-absorption modulator type optical modulator integrated laser diode having a different device structure. An optical modulator element and a laser diode element are arranged on a semi-insulating semiconductor board and an anode of each device is electrically isolated from its cathode so that driving by a single power source can be achieved.

Semiconductor element mounting board and optical transmission module


A computer-implemented method includes receiving a request for a web page; retrieving information associated with the web page, wherein the information comprises a link and one or more link placeholders associated with the link; determining context information associated with the computing device; generating, based on the context information, behavior information for the link; and populating at least one of the one or more link placeholders with the behavior information.

Generally, a web page is based on Hypertext Markup Language ("HTML") code, which provides the ability to display "links" in the web page. The term "link" refers to a segment of text or a graphical item that serves as a cross-reference between parts of a web page and/or between web pages. For example, through selection of a link, a viewer of a web page is able to navigate to a new web page.
A link may be associated with "link information." Generally, link information refers to information associated with a link. For example, link information may include target address information. A target address generally refers to an address that is loaded by a web browser and/or an application, for example, upon selection of a link. Link information may also include information specifying a behavior for the link ("behavior information"), including, e.g., how the link is displayed, whether a selection of the link launches a new application and/or a new browser window, whether alert messages are generated upon selection of the link, and so forth. Link information may also include information specifying whether a web page associated with a link is opened in a new browser window or is opened through "tabbed browsing" in an existing browser window. Generally, tabbed browsing refers to a function of some web browsers that allows a user to view multiple web pages in a single browser window by loading a web page into a section (referred to as a "tab") of the browser window, rather than in multiple browser windows.
The behavior of a link may depend on a type of device in which the link is displayed. For example, the type of link behaviors supported by a mobile device may differ from the type of link behaviors support by a laptop computer. In an example, a mobile device may not support tabbed browsing. In such a mobile device, a displayed link cannot be associated with tabbed browsing behavior.

Generating behavior information for a link


A transmission belt that is to be mounted on pulleys, for transmitting a force between the pulleys. The transmission belt includes: at least one endless annular hoop; and plate-shaped elements that are held by the at least hoop such that the elements are arranged in a circumferential direction of the transmission belt. The elements, each of which has a pair of flank surfaces at its widthwise opposite ends, are to be clamped between conical surfaces of a substantially V-shaped annular groove of each of the pulleys. A flank angle of each of the flank surfaces is larger than a cone-generating angle of each of the conical surfaces. The elements include first and second elements, such that each of the first elements is to be in contact at the flank surfaces with the conical surfaces, and has a width larger than a width of each of the second elements.

There is known a transmission belt to be mounted on a plurality of pulleys, for transmitting a force between the plurality of pulleys, for example, in a belt-type continuously variable transmission. The transmission belt includes (a) an endless annular hoop and (b) a multiplicity of plate-shaped elements that are held by the hoop such that the elements are arranged in a circumferential direction of the transmission belt, wherein each of the elements has a pair of flank surfaces at its widthwise opposite ends, and is to be clamped between a pair of conical surfaces of each of the pulleys, with the flank surfaces being in frictional contact with the conical surfaces of the pulley, for performing the force transmission between the pulleys. A transmission belt, which is disclosed in JP-2008-223938A, is an example of such a transmission belt. In the disclosed transmission belt, a plurality of types of elements having flank surfaces with different flank angles, so that the elements have various degrees of rigidity, for thereby reducing vibration and noise generated when each of the elements enters (or bites into) an annular groove of the pulley.