A system for including a set of software components and an application programming interface for the logical networked representation and implementation of a system of smart devices within a home or small office environment. The software components include a control object component and a user interface component. The control object is a software object capable that carries out communication, accepts and issues control messages. The control object allows a user to remotely control or query a physical device for which the specific control object has responsibility and for maintaining logical and physical connections and control of the device. The user interface supplies a coherent and transparent interface across all network media and topologies. The user interface is capable of communicating with any control object that it identifies as being controllable in a given system. The communication may be across a logical network or may be local to a particular device. The user interface takes on the “personality” of the physical device for which the attendant control object carries a logical and physical mapping and control responsibility.

In the past, environments having networked equipment were typically limited to commercial and industrial sites. Such equipment was typically complex and highly sophisticated computer controlled machines that utilized networks to coordinate manufacturing activities. However, unlike the commercial and industrial setting, the typical consumer setting included devices and appliances that were fairly unsophisticated and that were operated by a user on a one-to-one basis. To achieve the consumer device's specific function, the user had to physically interact with the device by, for example, pressing buttons on a coffee maker to preset the machine to turn ON at a specific time.
Early advances in consumer devices and appliances were limited to remote control devices that, at first, were wired to the device, and later used infrared (IR) communications to transmit commands. However, the remote control devices still forced the user to interact with their respective devices on a one-to-one basis. In addition, most hand-held remote control devices were proprietary, which required separate remote controllers for each consumer product. Further advancements included universal remote controllers that were capable of being programmed to control many consumer devices. However, universal remote control devices were limited in that they were typically not able to control devices other than entertainment equipment.
Lately, consumer appliances have become more sophisticated because of the decreased size and costs of computer-related components such as microprocessors. These new consumer devices include the so-called "smart" devices that, like the sophisticated commercial and industrial machinery, are capable of executing predefined tasks and capable of communicating over a network infrastructures. An example of a system is smart devices is CEBus®, which may utilize household powerlines to transport messages between compliant CEBus® household devices. CEBus® is implemented per EIA IS-60 for home automation. The CEBus® system utilizes controller and receptacle units to communicate and control CEBus® compliant devices. The receptacle units are connected to standard household receptacles and the command and control information is communicated over the powerlines. However, the CEBus® system is a hardware solution, which requires that the compliant devices be designed to be compatible with the standard. As a further limitation, the CEBus® controller is a special-purpose device. While this is an advancement over prior solutions, this hardware solution limits the upgradablity of the network and the devices, and fails to provide a solution for non-compliant smart devices.
A further example of a system that is intended to network devices is Jini™, which is based on Java™. Jini™ uses Java Remote Method Invocation protocols to move code around a network of devices. Network services run on top of the Jini™ software architecture such that devices and applications may register with the network via a lookup service. However, the Jini™ system is limited in that each device must register with the lookup service for other devices to discover them. Thus, each device is not aware of other devices on the network without first pointing to the lookup service. Further, there is no general purpose user interface that enables a user to interact and control all of the networked devices from a remote location.
In view of the above, there is a need for a system whereby once isolated and non-remotely controllable and configurable devices are able to share information with and about each other, and are able to provide and accept common command, control, and communication though a network. There is also a need for a system that provides the networked devices with the ability to operate autonomously in the case of a network failure, or for localized control and operation. There is also a need for a system that provides users with a method to control devices with a general purpose device (e.g., personal computer), and that provides a method to communicate over open network infrastructures. The present invention provides a solution to address the limitations and shortcomings of the prior art.

Control object and user interface for controlling networked devices


The invention includes a semiconductor construction having a wire bonding region associated with a metal-containing layer, and having radiation-imageable material over the metal-containing layer. The radiation-imageable material can be configured as a multi-level pattern having a first topographical region with a first elevational height and a second topographical region with a second elevational height above the first elevational height. The second topographical region can be laterally displaced from the bonding region by at least a lateral width of the first topographical region, with said lateral width being at least about 10 microns. Additionally, or alternatively, the elevational height of the second topographical region can be at least about 2 microns above the elevational height of the first topographical region. The invention also includes a method of forming wire bonds for semiconductor constructions in which a multi-level pattern is photolithographically formed in a radiation-imageable material (such as, for example, polyimide).

Wire bonds are commonly utilized for connecting integrated circuitry associated with semiconductor constructions to other circuitry external of the constructions. An exemplary prior art method of forming a wire bond for a semiconductor construction is described with reference to FIGS. 1-7.
Referring initially to FIG. 1, a semiconductor construction 10 is illustrated at a preliminary processing stage of the prior art method. Construction 10 comprises a semiconductor substrate 12 which includes a base 14, a conductive material 16 supported by the base, and an electrically insulative cap 18 over the conductive material.
Base 14 can comprise a semiconductor material, such as, for example, a monocrystalline silicon wafer having numerous integrated circuit devices (not shown) supported thereby.
Conductive material 16 can correspond to, for example, a metal-containing level formed over and in electrical connection with various integrated circuit devices associated with base 14. Conductive material 16 can, for example, correspond to the so-called level 1, level 2, level 3, level 4, level 5, etc. metal layers commonly associated with semiconductor constructions. In particular aspects, conductive material 16 can comprise, consist essentially of, or consist of aluminum and/or copper. In such aspects, material 16 can, for example, comprise one or more alloys comprising one or both of copper and aluminum. Conductive material 16 has an upper surface 17.
Insulative material 18 is a protective material formed over the uppermost surface of conductive material 16, and can comprise any suitable insulative material or combination of insulative materials. In particular aspects, insulative material 18 will comprise, consist essentially of, or consist of one or both of silicon nitride and silicon dioxide, For instance, insulative material 18 can comprise a homogeneous layer of silicon dioxide or silicon nitride. Alternatively, material 18 can comprise a multi-layer stack, with particular layers being silicon dioxide and other layers being silicon nitride. Insulative material 18 has an uppermost surface 19 (the label "19" is shown in FIG. 4).
Substrate 12 can be referred to as a semiconductor substrate. To aid in interpretation of the claims that follow, the terms "semiconductive substrate" and "semiconductor substrate" are defined to mean any construction comprising semiconductive material, including, but not limited to, bulk semiconductive materials such as a semiconductive wafer (either alone or in assemblies comprising other materials thereon), and semiconductive material layers (either alone or in assemblies comprising other materials). The term "substrate" refers to any supporting structure, including, but not limited to, the semiconductive substrates described above. In accordance with the definition of substrate provided above, all of the structures 14, 16 and 18 can be together considered to correspond to a semiconductor substrate. Alternatively, structure 14 can be considered to alone correspond to a semiconductor substrate. As yet another alternative, structure 14 and structure 16 can be considered to correspond to a semiconductor substrate, with structure 18 being considered to be formed over such semiconductor substrate.
A patterned radiation-imageable material 20 is formed over the uppermost surface 19 of insulative material 18. Radiation-imageable material 20 can correspond to photoresist, and can be photolithographically formed into the shown pattern. The photolithographic patterning of material 20 would typically comprise exposure of material 20 to a pattern of radiation, with such radiation being suitable to render exposed portions of material 20 either more or less soluble than non-exposed portions of the material in a solvent. The material 20 is then exposed to the solvent to remove the more soluble portions and thus form the shown pattern. The shown patterned material 20 has an opening 22 extending therethrough to an upper surface of insulative material 18.
Referring next to FIG. 2, the opening 22 is extended through insulative material 18 to an upper surface of conductive material 16 with an appropriate etch.
Referring next to FIG. 3, patterned material 20 (FIG. 2) is removed.
Referring next to FIG. 4, a patterned protective material 30 is formed over insulative material 18. In some aspects, insulative material 18 can be considered a passivation layer (or layers), and material 30 can be considered a protective cap formed over the passivation layer (or layers). Material 30 can correspond to a photolithographically patterned radiation-imageable material, and in some aspects comprise, consist essentially of, or consist of polyimide. The patterned material 30 defines an opening 32 extending to an upper surface of material 18. Opening 32 is wider than the opening 22 described previously (FIG. 3).
FIG. 5 is a top view of the FIG. 4 construction, and shows the wide opening 32 entirely surrounding narrow opening 22. Openings 22 and 32 are shown having rectangular peripheries, but it is to be understood that the openings can be formed in any suitable configuration, including, for example, configurations with curved peripheries.
Referring next to FIG. 6, a tool 40 is utilized to bond a wire 42 to conductive material 16. The wire 42 is retained within a capillary 44 of tool 40, and accordingly the tool can be referred to as a bonder capillary tool. The bonding of wire 42 to material 16 is accomplished by contacting the wire onto the upper surface 17 of material 16 with typically known in the art ultrasonic, thermal or other energy to weld the wire 42 and material 16 together. The region of material 16 where wire 42 connects with the material can be referred to as a wire bonding region of the material 16, and is indicated by the label 43 in FIG. 6. More specifically, a wire bonding region of material 16 is the portion of material 16 that ultimately directly contacts wire 42 in forming a bond to the wire. The wire can comprise, for example, gold.
Tool 40 has a shown lower portion proximate material 16, and such portion has a lateral width 50. Such lateral width which can be referred to as the lateral footprint of tool 40 proximate a bonding region. In the shown aspect of the prior art, such lateral footprint problematically extends over masking material 30 and accordingly the masking material is smashed downwardly by tool 40 during the bonding of wire 42 to material 16. The smashing of material 30 leads to smashing of the underlying material 18, which forms a damage regions 52 of material 18. The damage region can problematically contain cracks or other undesired structural flaws.
Referring to FIG. 7, tool 40 is removed to leave wire 42 bonded to the wire bonding region of material 16. Unfortunately, the damaged regions of materials 30 and 18 remain after removal of the tool. Accordingly, it is desired to develop new methods for forming wire bonds which alleviate formation of damaged regions within materials 30 and 18.
Although only one wire bond is shown formed in the diagrams of FIGS. 1-7, it should be understood that a large number of wire bonds is typically formed over a single substrate, with each wire bond being fabricated with the methodology of FIGS. 1-7.
A continuing goal of semiconductor processing is to reduce photomasking steps, in that each photomasking step carries with it a risk of mask misalignment and defect creation. Accordingly, it is desired to develop methodologies of forming wire bonds which reduce photomasking steps relative to the number of steps utilized in the processing sequence of FIGS. 1-7.
Although the invention was motivated, at least in part, by the problems discussed above, the invention is not limited to solutions of such problems.

Methods of forming wire bonds for semiconductor constructions


A peep sight device for a bowstring includes a pair of interengaged inner and outer sections, with a sight hole surrounded by a peripheral surface having a degree of taper such that substantially about 120.degree. of natural light is available to the sight hole on each of two opposite sides of the device. The sight hole is formed in the inner section, with an adjacent, interior groove diverting a portion of the bowstring away from the bowstring centerline, thereby enabling the sight hole to lie on the centerline.

In the field of archery, it is conventional practice to mount a peep sight on the bowstring, and to keep the sight in proper rotational alignment for accurate sighting of the target by the archer. It is difficult, however, to achieve precise rotational alignment of the peep sight, particularly with sights which are mounted between strands of the bowstring, and to maintain the alignment during a full draw.
The object of this invention is to provide an improved peep sight which is easy to install and precisely locatable on the bowstring. In addition, the unique design of the peep sight provides increased natural light to both sides of the sight hole.
In the exemplary embodiment, the peep sight is of two piece construction, with an inner piece or section being slidably receivable in a sideways direction within an outer piece or section. The sections are adapted to come together from opposite sides of the bowstring in a direction perpendicular to the latter, and in such a way as to sandwich the bowstring between the two sections. More specifically, the bowstring is moved laterally out of alignment with the longitudinal centerline of the bowstring, but retained within the peep sight, so as to insure alignment of the sight hole with the bowstring centerline. This is achieved by providing a groove on an interior side of the inner section of the peep sight, so that the bowstring is captured in the groove as the inner section moves into interfitting engagement with the outer section. More specifically, the groove is located in the free end of a transversely oriented stem, ahead of the sight hole in the direction of assembly/installation. After the two sections are interengaged along the bowstring, a fastening pin is inserted into aligned holes in the inner and outer sections.
Another feature of the invention lies in the provision for increased natural light on both sides of the sight hole. This is achieved by a beveled surface surrounding the sight hole which is formed in the inner section. A larger opening in the outer section aligns with the sight hole when the two sections are assembled, with a beveled surface surrounding the larger opening forming a continuation of the beveled surface surrounding the sight hole. These surfaces are tapered at an angle of about thirty degrees in all directions from the center axis of the hole (relative to a plane through the center of the peep sight separating the peep sight into front and rear halves), thereby enabling about 120.degree. of natural light to reach the sight hole on each side thereof. The same configuration is provided on the opposite of the sight hole, thereby providing the same amount of light to the opposite side of the sight hole. As a result, better sighting is possible and low-light shooting time is extended.
The two piece assembly also facilitates quick and accurate installation without need for a bowpress, and since the peep sight is usable from either side, the amount of turning required in the event of bowstring twist is reduced.
In its broadest aspects, therefore, the present invention thus relates to a peep sight device for a bowstring comprising a pair of interengaged inner and outer sections and including a sight hole surrounded by beveled surfaces on both sides thereof having a degree of taper such that substantially about 120.degree. of natural light is available to the sight hole on each of two opposite sides of the device.
In another aspect, the invention relates to a peep sight device for a bowstring comprising a pair of interengaged inner and outer sections, the inner section having a sight hole formed therein and an adjacent groove for diverting a portion of the bowstring away from a longitudinal axis of the bowstring, thereby enabling the sight hole to lie on the longitudinal axis.
Other objects and advantages will become apparent from the detailed description which follows.


A cell culture medium and cryopreservation medium in which sodium chloride is replaced with an organic cation, preferably choline chloride in a concentration of at least 100 mM, resulting a residual sodium ion concentration less than about 50 mM. The cryopreservation solution is suitable for cryopreservation of unfertilized oocytes, with thawed oocytes demonstrating the ability to survive, fertilize, and for the resulting embryos to proceed to full term development.

The ability to cryopreserve mammalian oocytes with high success rates in an easily reproducible manner has not yet been achieved. Cryopreservation of mouse oocytes, for example, could be valuable for the long term storage of genetically important strains of mice and would serve as a model for oocyte cryopreservation in humans, commercial livestock, and endangered species. This technology could also help alleviate the legal, ethical, and religious problems associated with the storage of human embryos, as well as improve current embryo freezing protocols. Although offspring have been produced from frozen-thawed oocytes in the mouse, human, bovine, and rabbit (Whittingham, 1977; Al-Hasani et al., 1987; Fuku et al., 1992; Chen 1986; Van Uem et al., 1987), results have been variable and not sufficiently successful to make oocyte cryopreservation routine.
In addition to oocyte freezing, cryopreservation and transplantation of ovarian autografts have been somewhat successful in mice and sheep (Gosden et al., 1994; Gunasena et al., 1997). The ovarian tissues in these reports were frozen using simple embryo freezing protocols. Ovarian tissue survived freezing and was able to continue development after transplantation to the reproductive tract or kidney capsule, yielding growing follicles.
Current cryopreservation protocols have evolved from embryo freezing methods that produced offspring in mice, cows, and sheep (Whittingham et al., 1972; Willadsen et al., 1976, 1978). Cryopreservation procedures for mouse and other mammalian embryos are now relatively efficient, but these techniques cannot be used reliably for oocyte freezing. Studies cryopreserving mouse oocytes report very different survival and fertilization rates (Carroll et al., 1993; Carroll et al., 1990; Cohen et al., 1988; George et al., 1994; Glenister et al., 1990; Gook et al., 1993; Whittingham et al., 1977). The variability in the success of mouse oocyte freezing is most likely due to modifications in freezing protocol and/or the use of different cryoprotectants. Although these protocols differ, they rely on the same basic cryobiological principals.
A particularly important area that would benefit greatly from advances in cryopreservation technology is assisted human reproduction. In this field, one or both partners may have a fertility problem. In order to overcome these problems, eggs are harvested from the mother, and sperm from the father. The sperm are then used to fertilize the eggs, and one or more developing embryos are then replaced to the uterus of the mother. Typically, egg maturation is induced pharmaceutically prior to harvesting, and the sperm must be available to fertilize the eggs. Generally, only a very limited number of surgical harvesting procedures may be conducted on an individual, and the number of eggs replaced in the mother is limited in order to avoid multiple live births. Therefore, there has been a need to preserve harvested eggs, sperm, fertilized eggs, and embryos.
The ability to successfully and reproducibly cryopreserve oocytes, which is not yet possible, would allow women to have their eggs frozen until a time when they have found a suitable sperm donor (possibly a future husband) and then thaw their frozen eggs and fertilize them. The resulting embryos can be replaced into the patients' uterus, after it has been prepared to receive the embryo using hormones and known techniques, to allow for implantation, fetal development, and ultimately birth. Ethically, there are issues involved in preserving fertilized eggs or embryos, that are more significant than those associated with freezing unfertilized germ cells. Cryopreservation of oocytes would avoid the ethical concerns surrounding embryo freezing in humans and offer another option to couples with infertility problems. Storing oocytes early in life, such as when the mother is in her 20's and early 30's, when healthy eggs tend to be produced, and therefore when they are more apt to be fertilized and result in viable offspring, would greatly improve the chance of a pregnancy later in life, rather than the relatively poor pregnancy rates produced when participating in an in vitro fertilization (IVF) program at 40 years old or older.
Cryopreservation of oocytes, especially from humans, in a reproducible and efficient manner has been generally unsuccessful according to known techniques. It is noted, however, that offspring have been produced from frozen oocytes of several species, including humans. An improved cryopreservation medium would benefit oocyte storage and may also provide a more successful way of freezing embryos, thereby improving the possibilities for pregnancy.
Cryopreservation technology is applicable to other species besides humans. In commercial livestock, improvements in oocyte or embryo cryopreservation could greatly improve genetic management of populations and the number of offspring generated, resulting in a significant time savings and efficiency. In endangered animals, any improvement in embryo freezing or the development of a method to freeze eggs could lead to an increase in the number of offspring produced, enhancement of the genetic quality of the offspring, improvement of the population's genetic health, elimination of both the cost and risk of transporting live animals for reproductive purposes, and possibly even a delay or prevention from extinction of certain species.
Cryopreservation of oocytes from commercial livestock would provide a valuable means of saving and distributing important genetic material. Oocytes from animals exhibiting commercially valuable traits such as increased milk production could be stored.
Cryopreservation of oocytes from endangered species would provide an invaluable method of salvaging important genetic material. Provided that the eggs could be thawed, fertilized, and produce fertile offspring, and given that sperm is relatively easy to store, species could be stored indefinitely, virtually eliminating the risk of extinction. Cryopreserved oocytes could be easily transported globally, providing a source of genetic information to better aid in managing populations. In this manner, underrepresented genes from founder animals could be reintroduced into the population at any time, even 200 years from now. Because frozen oocytes (and sperm) can be distributed easily and cost effectively, the possibilities of improving genetic diversity and the overall health of a population are intriguing. Oocytes collected in the field can be frozen and infused into the captive population to improve its genetic health. With this technology in place frozen zoos can become a reality.
Biological cryopreservation systems are well known. They allow cells to be frozen, for example at -20.degree. C. or below, for extended periods, and resume normal cell activity after thawing. Typically, problems encountered include intracellular ice formation (IIF) and osmotic imbalances that result in cellular disruption. It is the prevention of IIF that many prior methods are directed.
Cryopreservation of cells involves dehydration, introduction of a cryoprotectant, and cooling to a low temperature, usually from -30.degree. C. to -80.degree. C., before plunging in LN.sub.2. The first objective is the removal of water from the cell, which when cooled below its melting point forms ice crystals that can damage intracellular organelles as well as the cell membrane (Mazur, 1977). The osmolality of the extracellular solution increases as the outside water freezes, causing the water to passively exit the cell. More ice forms at lower temperatures resulting in continued cellular dehydration. The next objective for freezing cells concerns the combining of any residual water left in the cell with a cryoprotectant, in order to form a glass-like structure when solidified, thereby preventing IIF. Because the melting point of water is reached both during freezing and thawing, IIF can occur at either time. Damage may therefore occur when the cell is exposed to elevated concentrations of electrolytes and/or IIF. IIF can be catalyzed by the presence of extracellular ice surrounding the cells (seeding) or heterogeneously by intracellular structures. In the presence of cryoprotectants however, IIF resulting from either seeding and/or heterogeneous ice nucleation does not occur (Toner et al., 1991), suggesting that IIF may not be a problem when freezing oocytes in the presence of a cryoprotectant. Therefore, electrolytes including sodium would appear to impart the majority of the damage during cryopreservation. Cellular disruption by sodium ions could alter the cell membrane and/or intracellular organelles. Lovelock (1954) hypothesized that cellular damage is caused by an increase in electrolyte concentration, causing destabilization of membranes. Mazur et al. (1974) further demonstrated that the cell surface is a major site of freezing damage. In his study, Mazur showed that the nonpermeating solute sucrose was as effective in protecting erythrocytes from freeze/thaw damage as the permeable cryoprotectant glycerol. The majority of damage to mouse oocytes during freezing may be caused by sodium ions, but we cannot rule out the possibility of IIF. Whether cellular demise is caused by exposure to elevated sodium ion concentrations during the freeze or thaw, or a problem with the transport of sodium ions across cell membranes remains to be elucidated. Sodium ions have a radius of about 0.95 .ANG., while lithium has a radius of 0.60 .ANG. and potassium has a radius of about 1.33 .ANG.. The majority of past cryopreservation studies that have focused on IIF, cryoprotectants, and altering freezing protocols have been unable to significantly improve oocyte freezing. Therefore, the type of cryoprotectant used, or the freezing protocol may already be adequate for oocyte freezing and IIF may not be the major problem presumed by the prior art.
Standard embryo cryopreservation techniques are remarkably similar. In general, embryos are exposed to a cryoprotectant (dimethyl sulfoxide (DMSO), 1,2-propanediol (PrOH), glycerol, ethylene glycol), diluted in a simple sodium-based salt solution for 5-15 min to allow uptake of the cryoprotectant. The embryos are then cooled quickly (-2.degree. C./min) to about 7.degree. C. at which point they are seeded, cooled slowly (-0.3.degree. C. to -0.5.degree. C./min) to about -30.degree. C. or below, and then plunged directly into liquid nitrogen (LN.sub.2). Embryos can also be rapidly frozen or vitrified, but only using very elevated cryopreservative concentrations (2M to 6M) that are toxic to cells when they are exposed for more than a few minutes. Following cryopreservation the embryos are thawed and cultured. Thawing procedures differ, but very little. Two basic concepts are involved in the thawing process, 1) removal of cryoprotectant and 2) rehydration of the embryo at a rate so as not to rupture the cell membrane, usually with the use of sucrose. These freezing and thawing procedures work relatively well for embryos, but do not allow the successful storage of oocytes. The exact reason(s) why embryos can survive and oocytes cannot are unknown, but membrane damage due to IIF, ion loading, and/or osmotic stress are suspect. Researchers have mainly focused on the problem of IIF and osmotic stress by modifying freezing procedures (slow vs. fast), thawing procedures (slow vs. fast), and methods for the removal of the cryoprotectants with little attention to the cryopreservation media formulation, which is presumed to be optimal.
Researchers have tried unsuccessfully to cryopreserve oocytes in a reproducible manner using one or more cryoprotectants. In general, PrOH and DMSO are the most common cryoprotectants used today for the cryopreservation of oocytes and embryos. Although all of the pregnancies in the human have resulted from oocytes frozen in DMSO, PrOH is the cryoprotectant of choice because of its greater permeability, reduced toxicity, and improved success in storing human embryos (Gook et al., 1995; Imoedemhe and Sigue, 1992; Lassalle et al., 1985; Mandelbaum et al., 1987; Testart et al., 1986). Mouse oocytes have been frozen using PrOH, but with poor overall results (Gook et al., 1993; Todorow et al., 1989). Todorow et al. (1989) reported survival and fertilization rates of 63% and 27%, respectively. When PrOH was used in combination with DMSO, survival and fertilization rates increased to 87% and 42%, respectively. Numerous studies have reported survival and fertilization rates of up to 79% and 50%, respectively, for mouse oocytes cryopreserved using DMSO (Glenister et al., 1987; George and Johnson, 1993; Carroll et al., 1989; Todorow et al., 1989; Schroeder et al., 1990; Aigner et al., 1992; Bouquet et al., 1992). Although one laboratory has reported survival rates of up to 91%, fertilization up to 78%, and blastocyst formation as high as 54.4% (Carroll et al., 1993), it remains to be seen whether their results can be easily reproduced.

Cryopreservation and cell culture medium comprising less than 50 mM sodium ions and greater than 100 mM choline salt


A zoom finder drive mechanism, driven only when taking a picture. When a lens barrel is in a collapsed position, the cam-follower pin of a rotary barrel does not abut at an end of a protrusion of a finder cam barrel such that a first and second movable finder lenses are idle. When taking a picture, the cam-follower of the rotary barrel is abutted at the end of the protrusion of the finder cam barrel such that the first and second movable finder lenses are driven.

The invention relates to zoom finder drive mechanisms, and in particular to zoom finder drive mechanisms driven by movement of a zoom lens barrel.
Operations of a zoom finder in a camera are performed and moved by an optical photography system. In a conventional driving method for a zoom finder, rotary lens barrel driving elements are used as a direct driving source. A zoom motor driving the zoom lens barrel and speed-reduction gears thereof are also used as direct driving sources.
U.S. publication No. 2004-0042092 discloses a method for driving zoom finders. When a lens barrel is retracted or housed therein, the zoom finder is idle. When taking a picture, only the lens rotary elements in the zoom region and the zoom finder are connected by gears.
The zoom finder, driven in the zoom region when taking a picture, can reduce the moving stroke of the driving elements or prevent occurrence of needless loading in a region other than the zoom region. In U.S. publication No. 2004-0042092, however, since the lens rotary elements in the zoom region and the zoom finder are connected by gears, the gears are necessary components, increasing complexity and reducing available space therein. Thus, costs are increased.


An article of waterproof footwear and its method of manufacture are disclosed. The article of waterproof footwear has an upper containing a water-impermeable layer. The distal edge of the upper has a polymeric binding covering the distal edge and the inside and outside surfaces of the upper adjacent to the distal edge. The polymeric binding seals to a waterproof insole thereby effectively waterproofing the upper of the footwear. An outer sole is attached to the upper thereby forming the waterproof footwear.. The method of manufacture allows production of a waterproof upper and the attachment to the upper of an outer sole, thereby forming the waterproof footwear of the invention.

In order to produce waterproof footwear, footwear is manufactured totally from rubber or another polymer through a dipping or molding technique, such that there is no separation between the upper and the outer sole, and therefore the footwear is water-impermeable. However, waterproof footwear produced through a dipping or molding technique has poor fit, is heavy in weight, and is impermeable to water vapor. Additionally, waterproof footwear produced through a dipping or molding technique is not readily adaptable to various footwear styles.
In recent times, footwear was manufactured whose body or upper part consists of water-impermeable, and preferably, water vapor-permeable material. There have been problems in providing a waterproof connection between the upper and an outer sole of the footwear while still maintaining the good fit, lightweight, and water vapor-permeable quality of the upper material.
In an attempt to produce waterproof footwear whose upper part consists of water-impermeable and water vapor-permeable material, a cement process of lasting footwear was developed. In this process, an upper of a shoe is cemented to an insole. To this unit, a sole is applied which may be an intermediate sole or an outer sole. It is a problem to have truly waterproof footwear with cement-lasted footwear, even if a water-impermeable outer sole and a water-impermeable upper layer are utilized in the construction. The weak point in the production of waterproof cement-lasted footwear is the formation of a durable waterproof seal between the insole and the upper since the lasting cement does not initially seal between the upper and the insole and may become brittle and more water-permeable due to bending stresses during use of the footwear. The formation of a durable waterproof seal between the insole and the upper is hampered in the lasting process since pleats form in the edges of the upper materials because the straight edges of the upper are forced to lay flat against the curved insole. These upper wrinkles are three-dimensional in nature and therefore provide easy routes for water entry into the upper of the footwear.
To overcome problems with the cement process of lasting footwear, a waterproof insert method was developed wherein a unit of a footwear upper and a cemented insole is lined with a sock-like insert of a water impermeable, water vapor-permeable material, as taught in U.S. Pat. No. 4,599,810, to Sacre. To obtain an insert that is waterproof and of the desired shape, the insert must first be sewn and then hot-welded at the sewn seams to form a waterproof article. This waterproof insert method does not allow the insole to be directly attached to a shoe last in a single step as in the traditional cement process of lasting footwear. An additional lasting step is usually required by this method, making this method more complicated and expensive for most shoe manufacturers.
An injection molding process for the soling of footwear, as taught in U.S. Pat. No. 4,899,465 to Bleimhofer, et al., has been used to produce waterproof footwear. A polyurethane outer sole is molded by machine to the sole region of an upper. It is necessary to have an injection mold that is applied to the upper from both sides. This injection mold is relatively expensive. Due to the high mold costs, the shoe manufacturers are restricted to very few sole configurations. Besides, it is not possible to achieve footwear having an elegant appearance with such molded-on outer soles.
The instant invention is directed to an improvement of the cement process of lasting footwear in such a manner that the connection between the upper of the footwear and insole region is waterproof in a reliable manner, while permitting any kind of outer soles to be employed.


A method for fabricating a chip size package is provided. The method includes the step of forming a laminated substrate which consists of a dielectric layer and a highly conductive layer disposed thereon. Holes are drilled into the dielectric layer. A desired pattern is applied to the conductive layer. A chip structure is formed which consists of a silicon die and an insulating layer disposed thereon. Gold bumps are applied to the top surface of the bonding pads. The laminated substrate is bonded to the chip structure via the holes and gold bumps. A solder mask is applied over the top surface of the conductive layer of the laminated substrate so as to form selective solder areas. Finally, solder balls are attached to the selective solder areas.

This invention relates generally to semiconductor integrated circuit packages and more particularly, it relates to an improved method for fabricating a chip size package which can accommodate high or low pin-count dies with a minimum amount of substrate real estate being required.
In recent years, there has been an increased interest in the packaging industry in an area generally known as a die size solution. The reason for the search for new alternatives is that the package and interconnection industry has moved away from the use of pins as connectors for electronic packaging due to the high cost of fabrication, the high failure rate in the connections, and the limitation on the density of input/output pins. Consequently, there exists a need to eliminate the classic package performance problems and to reduce the size, weight, and cost of standard packages. As a result, there has been developed new die size packaging alternatives which are classified into three main groups: (1) Bare die or Known Good Die (KGD), (2) flip chip, and (3) chip size package (CSP).
With respect to the chip size package (CSP) alternative, the package typically occupies in the range of 1 to 1.2 times the area of the chip in the package. The primary advantage of the CSP is that they can accept a die design with peripheral bond pads and route the signal to a ball on the surface to make it suitable for surface mounting. The rerouting approach on the surface of the die reduces the handling and mounting issues over the other two mentioned alternatives and facilitates in the testing processes of the product. However, the conventional versions of the CSP alternatives generally use a wire bond technique to connect the bond pads to an interposer surface for holding the solder balls, thereby introducing increased electrical resistance and conductance.
It would therefore be desirable to provide an improved method for fabricating a chip size package which is relatively simple in construction and easy to assemble. It would be also expedient to provide a chip size package which eliminates the need for TAB (tape-automated bonding) or wire bonding in order to reduce electrical resistance and conductance.

Process for manufacturing a multi layer bumped semiconductor device


A method for determining file fragmentation and performing subsequent defragmentation, including measuring a file fragmentation factor, measuring a file extent distance factor, measuring a file fragmentation level, selecting a file stored as multiple extents on at least one disk, selecting a free space extent on the disk whose length most closely matches the length of the file, calculating a fragmentation level of the disk from a fragmentation factor and a file extent distance factor of free space extents or allocated space extents of the disk, calculating the effect that moving the file to the free space extent would have on the fragmentation level, and moving the file to the free space extent provided that doing so causes the fragmentation level to satisfy a criterion.

As files are created and deleted on a Direct Access Storage Device, hereinafter referred to as a "disk," the files and the free space on the disk typically become fragmented. Fragmentation occurs when an individual file is not stored to a contiguous region of the disk, but rather is broken up into two or more pieces, or "extents," which are stored to separate regions of the disk. When a program tries to access a fragmented file, the disk read/write head must often move between multiple disk regions in order to access each extent of the file. As more files become fragmented, and as individual files are fragmented into more and more extents, more read/write head movements are required for file access, resulting in reduced performance.
Defragmentation may be employed to increase the amount of contiguous free space on the disk and reduce the number of file extents, thereby reducing head movement and file access time. In one common method for defragmenting a disk, all of the files on a disk are copied to another storage device. After the files are copied, they are then deleted from the disk. The files are then copied back onto the disk such that each file resides in contiguous regions of the disk space. This method is effective, but requires a second, spare storage device to be used as the intermediate storage medium.
In another method, file extents are moved to other disk regions to create a contiguous area of free space large enough to receive a single file. This process typically continues until all files are stored contiguously, with each file stored as a single extent, and with all free space being contiguous. While this method does not require a separate, intermediate storage medium, it is relatively slow due to the extensive data movement required.
When and how to defragment a disk often first depends on measuring the level of fragmentation of the disk. In one method, the level of fragmentation is expressed as a function of the number of free space fragments on the disk relative to the total free space. However, this method does not reflect qualitative measures of fragmentation. For example, while two files having the same number of extents may have the same quantitative measure of fragmentation, a qualitative fragmentation measure might reveal a significant difference between the two files, such as where the extents of one lie along the same disk track and on the same disk platter, or along the same cylinder on multiple disk platters, requiring little head movement, while the extents of the other lie in multiple disk tracks and cylinders on multiple disk platters, requiring greater head movement.
A qualitative fragmentation measure would therefore be advantageous, as would a defragmentation method that incorporates such a measure.

Measuring fragmentation on direct access storage devices and defragmentation thereof


A cementitious slurry containing blast furnace slag and a surfactant is utilized to displace an oil based drilling fluid without causing contamination by the blast furnace slag cement. In a specific embodiment, a cementitious slurry comprising blast furnace slag, a surfactant and water is used in the cementing of an annulus surrounding a casing or liner wherein a synthetic oil-containing drilling fluid is present.

This invention relates to drilling and cementing boreholes.
The drilling of boreholes is generally carried out using a rotary drilling process. The rotary drilling of a borehole is accomplished by rotating a drill string having a drill pipe and a drill bit at its lower end. Weight is applied to the drill bit while rotating to create a borehole into the earth. The drill string is hollow and sections are added to the drill string to increase its length as the borehole is deepened. This rotary drilling process creates significant amounts of friction which produces heat along with fragments of the strata being penetrated. The fragments of the strata must be removed from the borehole and the drill bit must be cooled to extend its useful life. Both of these necessities are accomplished by the circulation of a fluid down through the drill string and up to the surface through the annulus between the drill string and the wall of the borehole.
Once the borehole has been drilled to the desired depth, it may be desirable to isolate the separate areas, zones or formations transversed by the borehole. For extraction of fluids from formations, and for protection against borehole collapse, one or more conduits (casings) must be inserted into the borehole extending from the surface downward. Thereafter, liners may be hung inside the casings.
At this point it becomes necessary to fill the annulus between the casing and the borehole wall or between the liner and casing with a cementitious material which will seal the annulus to inhibit communication between various formations penetrated by the wellbore and which will provide structural support for the casing or liner. This is commonly referred to as primary cementing.
Generally, the borehole, into which the casing or liner to be cemented is installed, is filled with drilling mud. Therein lie several problems. Conventional Portland cement and conventional drilling muds are incompatible. Thus, as the cement is forced down the casing or liner and up into the annulus it is commingled with the drilling mud at any interface between the mud and the cement. The resulting mixture generally thickens or becomes a gel and does not set up into a strong cement. In addition, the gel strength and viscosity become uncontrollable and the mixture may either become too viscous to pump or may get thinner. If the mixture gets thinner it can allow solids to settle downward in the annulus where they may bridge and restrict the passage of the cement slurry. In either event, the situation is unsatisfactory.
The industry has developed a complex system to attempt to circumvent this problem. A device generally known as a landing collar is attached to the bottom of the casing or liner being cemented. The landing collar has an annular shoulder projecting inwardly. A first wiper plug with a diaphragm which can be ruptured is introduced into the casing or liner followed by a spacer fluid, the cementitious slurry and finally, a solid wiper plug. Displacement fluid then pushes the solid wiper plug downward thus displacing the spacer fluid and the cementitious slurry through the ruptured first wiper plug, out of the casing or liner and into an area at the bottom thereof where it changes direction and flows upwardly into the annulus. When the second wiper plug reaches the area of the landing collar it is stopped by the first wiper plug which is resting on the shoulder.
The spacer fluid, however, can create its own set of problems because it mixes somewhat with both the drilling mud ahead of it and the cementitious slurry behind it. Usually, the most damaging is the contamination of the drilling mud. The spacer fluid remains a fluid in admixture with the drilling fluid and thus cannot be easily separated therefrom. In many instances this necessitates the disposing of the thus-contaminated drilling fluid. This is an economic problem with any drilling fluid but with more expensive, less environmentally friendly oil based fluids, it is a major obstacle to successful operation. In addition, the mixture of spacer fluid and drilling fluid usually results in a more viscous material thus exacerbating the pumping problems. This can even result in a total failure of the cement job because a good seal at the top of the liner may not be obtained and this is where any cement weakened by drilling fluid contamination most likely will be.
It has been found by Hale and Cowan, U.S. Pat. No. 5,058,679 (Oct. 22, 1991) that blast furnace slag is compatible with drilling fluids, and that excellent cementitious compositions can be produced by combining blast furnace slag and drilling fluids. Even blast furnace slag, however, can contaminate oil based drilling muds when the two are mixed. Parceveaux et al., U.S. Pat. No. 5,101,902 (Apr. 7, 1992) broadly discloses spacers between mud and cement and discloses adjusting the density with weighting agents including barite and fillers.
Oil based drilling fluids are disclosed in Shell Internationale patent GB 2,223,255A (Apr. 4, 1990). Cowan, U.S. Pat. No. 5,016,711 (May 21, 1991) broadly discloses improved cement adhesion through the use of a surfactant.

Oil mud displacement with blast furnace slag/surfactant


The present invention pertains to a semiconductor light-receiving element and a method for manufacturing the same, enabling operation in a wide wavelength bandwidth and achieving fast response and high response efficiency. A PIN type photodiode made by sequentially layering on top of the substrate a Si layer of a first conductivity type, a non-doped Ge layer and a Ge layer of a second conductivity type that is the opposite type of the first conductivity type and a Ge current-blocking mechanism is provided in at least part of the periphery of the PIN type photodiode.

As data transmission amounts between server and CPUs have increased in recent years, the handling of data through transmission of electrical signals using conventional Cu wires are approaching the limit. In order to resolve this bottleneck, optical interconnection, that is to say, data transmission through optical signals is required. Furthermore, from the points of view of reduction in power consumption and reduction in the area of devices, elements formed by integrating optical components such as optical transmitters, optical modulators and optical receivers that are required for optical transmission/reception on a Si substrate become necessary.
Meanwhile, optical components integrated on a Si substrate are connected through optical fibers and, therefore, it is preferable to use the wavelength of 1.55 μm, with which loss is low in optical fibers, as a transmission wavelength band. Therefore, it is preferable to use Ge, which has an absorption end in the vicinity of 1.55 μm, as an absorbing layer of a photodetector to be used for optical transmission with a band of the wavelength of 1.55 μm.
Meanwhile, wavelength division multiplexing (WDM) transmission is required as the data transmission amount increases. In order to implement WDM transmission, Ge photodetectors having a high response sensitivity in a broad region of which the wavelengths are longer than 1.55 μm are necessary.
In general, when the temperature is cooled to room temperature while Ge is growing on a Si substrate, the Ge epitaxial layer is subject to tensile strain in the directions within the plane of the substrate due to the difference between the coefficients of thermal expansion of Si and Ge. As a result, the absorption end of Ge on the Si substrate has a longer wavelength as compared to Ge layers in a bulk state, as has been reported (see Non-Patent Literature 1). This works advantageously from the point of view of expansion of the wavelength band of the photodetectors.
Meanwhile, the element capacitance of photodetectors is required to be lowered from the point of view of increase in the high-speed response properties and, thus, it is necessary to make the element area smaller (make the element width narrower). In addition, from the point of view of increase in the response sensitivity properties, it is required to prevent photocarriers generated in the depletion layer from recombining or becoming trapped by lattice defects while drifting. In order to do so, it is likewise necessary to make the element area smaller (make the element width narrower).
Thus, it has been reported that a photodetector having Ge on a Si substrate as an absorbing layer is processed as a mesa element having a width of approximately of several μm (see Non-Patent Literature 2). FIG. 20 is a schematic cross sectional diagram illustrating a conventional photodetector having Ge as an absorbing layer where the photodetector is formed using an SOI substrate. A Si layer 83 provided on top of a Si substrate 81 with a BOX layer 82 in between is processed to form a p type Si mesa portion 84 and p type Si slab portions 85 on the two sides thereof. At this time, a waveguide in stripe form is formed so as to be connected to the p type Si mesa portion 84 through a tapered portion, though this is not shown.
A non-doped Ge layer is formed on top of this p type Si mesa portion 84 through selective growth and n type impurities are implanted into the surface thereof so as to provide an n++ type Ge contact layer 87, where the portion into which the impurities have not been introduced is an i type Ge light absorbing layer 86. Meanwhile, p type impurities are implanted into portions of the p type Si slab portions 85 so as to form p++ type Si contact portions 88.
Next, an oxide film 89 that becomes the upper clad layer of the waveguide in stripe form is formed, plugs 90 and 91 are formed and, then, an n side electrode 92 and p side electrodes 93 are formed. Light that propagates through the waveguide in stripe form passes through the p type Si mesa portion so as to reach, and be absorbed by, the i type Ge light absorbing layer 86 through evanescent coupling.

Semiconductor light-receiving element and method for manufacturing same


A method for producing nickel composite hydroxide particles may include: a first step of producing nuclei including primary particles by controlling the pH of an aqueous solution for nucleation, the aqueous solution for nucleation containing a metal compound having an atomic ratio of metals corresponding to an atomic ratio of metals in the nickel composite hydroxide particles and substantially not containing a metal complex ion-forming agent; and a second step of forming, on an outer surface of each of the nuclei, an outer shell portion including platy primary particles larger than primary particles of the nuclei by controlling the pH of an aqueous solution for particle growth containing the nuclei obtained in the nucleation step.

In recent years, along with the spread of mobile electronic devices such as mobile phones and notebook-sized personal computers, development of smaller and lighter nonaqueous electrolyte secondary batteries having a high energy density has been strongly demanded. Development of high power secondary batteries as batteries for electric automobiles typified by hybrid automobiles has also been strongly demanded. The secondary batteries that meet such demands are exemplified by lithium ion secondary batteries. Lithium ion secondary batteries include a negative electrode, a positive electrode, an electrolytic solution and the like, in which a material into and from which lithium can be inserted and desorbed has been used as an active material for the negative and positive electrodes.
Research and development of the lithium ion secondary batteries have been extensively performed at present, and in particular, lithium ion secondary batteries in which a layer or a spinel type lithium metal composite oxide is used as a positive electrode material can achieve a voltage as high as 4 V; therefore, practical applications thereof as batteries having a high energy density have been accelerated.
As positive electrode materials for such lithium ion secondary batteries, lithium composite oxides such as lithium cobalt composite oxide (LiCoO2) which can be comparatively easily synthesized, lithium nickel composite oxide (LiNiO2) in which nickel less expensive than cobalt is used, lithium nickel cobalt manganese composite oxide (LiNi1/3Co1/3Mn1/3O2), and lithium manganese composite oxide (LiMn2O4) in which manganese is used have been hitherto proposed.
In order to achieve favorable performances of positive electrodes (superior cycle characteristics, low resistance and high power), positive electrode materials are required to be formed with particles having a uniform and appropriate particle diameter. The grounds for such requirements are that use of a material having a large particle diameter and a small specific surface area leads to failure in reserving a sufficient area for reaction with the electrolytic solution, thereby resulting in an increase of the reaction resistance and failure in obtaining a battery having a high power, and that use of a material having a broad particle size distribution leads to lowering of the battery capacity, thereby resulting in defects such as an increase of the reaction resistance. The reason why the battery capacity is reduced is that nonuniformity of a voltage applied to the particles in the electrode causes fine particles to selectively deteriorate due to repetition of charge and discharge.
Additionally, when increase in output of batteries is intended, shortening a distance of transfer of lithium ions between the positive electrode and the negative electrode is effective. Therefore, thin positive electrode plates have been desired, and thus cathode active material particles having a small particle diameter are useful also to this end.
When further increase in output of batteries is intended, cathode active material particles having a small particle diameter and a large specific surface area are useful. For example, by reducing smoothness of surfaces of the particles, employing a porous structure for the particles themselves, or the like, the specific surface area can be increased even though the particle diameters of the particles are the same. This leads to increasing an area for reaction between the particles and the electrolytic solution and enables to enhance reactivity therebetween. Increase in output of batteries can therefore be expected.
Additionally, it is necessary to produce particles so as to have a small and uniform particle diameter and a large specific surface area also in the case of the above lithium nickel composite oxide in order to improve performances of the positive electrode material.
Patent Literature 1 discloses a lithium composite oxide in the form of particles having a particle size distribution exhibited on the particle size distribution curve thereof in which: an average particle diameter D50 which means that the particle diameter of the particles with the accumulation frequency of 50% is 3 to 15 μm; the minimum particle diameter is 0.5 μm or more; and the maximum particle diameter is 50 μm or less, and with respect to the relationship between D10 which means that the particle diameter of the particles with the accumulation frequency of 10% and D90 which means that the particle diameter of the particles with the accumulation frequency of 90%, D10/D50 is 0.60 to 0.90, whereas D10/D90 is 0.30 to 0.70. The Literature also discloses that the lithium composite oxide has high fillability, favorable charge and discharge capacity characteristics and high output characteristics, and is less likely to deteriorate even under conditions with a significant charging and discharging load; therefore, use of this lithium composite oxide can provide a lithium ion nonaqueous electrolytic solution secondary battery having excellent output characteristics with small deterioration of the cycle characteristics.
Patent Literature 2 discloses a cathode active material for a nonaqueous electrolytic solution secondary battery having at least a lithium-transition metal composite oxide with a layer structure in which the cathode active material for a nonaqueous electrolytic solution secondary battery includes the lithium-transition metal composite oxide formed with hollow particles having an outer shell portion outside and a space portion inside the outer shell portion. The Literature also discloses that the cathode active material for a nonaqueous electrolytic solution secondary battery has excellent battery characteristics such as cycle characteristics, output characteristics and thermal stability and is used suitably for lithium ion secondary batteries and the like.
Since the lithium composite oxide particles disclosed in Patent Literature 1 are particles whose minimum particle diameter is 0.5 μm or more and whose maximum particle diameter is 50 μm or less with respect to the average particle diameter of 3 to 15 μm, the lithium composite oxide particles include fine particles and coarse particles. The particle size distribution defined by the D10/D50 and D10/D90 described above therefore does not suggest a narrow range of the particle diameter distribution. Consequently, the lithium composite oxide disclosed in Patent Literature 1 does not correspond to particles having uniform particle diameters, and thus improvement of the performances of the positive electrode material is not expected even when such a lithium composite oxide is employed, indicating difficulty in obtaining a lithium ion nonaqueous electrolytic solution secondary battery having sufficient performances.
Since the cathode active material for a nonaqueous electrolytic solution secondary battery disclosed in Patent Literature 2 includes hollow particles, increase of the specific surface area can be expected comparing with solid particles; therefore, improvement in reactivity between the particles and the electrolytic solution can be expected because of the increase of the specific surface area. Patent Literature 2 however does not describe a particle diameter and a particle size distribution of the cathode active material for a nonaqueous electrolytic solution secondary battery, considering it to be equal in terms of quality to a conventional cathode active material. Therefore, there is a high possibility that selective deterioration of fine particles occurs due to nonuniformly applying a voltage in an electrode, resulting in a reduction in a battery capacity.

Nickel composite hydroxide particles and nonaqueous electrolyte secondary battery


A multi-terminal block for electric devices, comprising two superimposed connecting levels for conductors, which can be used independent of each other, the multi-terminal block comprising: a first terminal block comprising a cage screw type conductor connector which comprises a cage and a screw for use securing the a first conductor, wherein the first terminal block cooperates with a terminal; and a second terminal block comprising a quick connecting type conductor connector, wherein the quick connecting type conductor connector comprises a strip spring for securing a second conductor, wherein the screw does not extend through the second terminal block and wherein the second terminal block is partially disposed within the cage such that the second terminal block cooperates with the terminal.

1. Field
The present disclosure relates generally to a multi-terminal block for electric devices, having two superimposed connecting levels for conductors, which can be used independent of each other. In particular, the double terminal connection of the present disclosure is used in MiniCircuit Breakers (MCB) and Residual Current Circuit Breakers (RCCB), wherein the terminal has two independent connections, i.e., a cage screw type and a second pin plug-in. The second terminal is preferably designed to connect a pin busbar without a screw, instead utilizes a metal sheet spring, with a shape that resembles the number "5". The plug-in is surrounded by the terminal cage belonging to the screw type connection. In order to access to the bladespring the pin busbar goes through the window of the terminal.
2. Discussion of the Background Art
Conventional double connection terminals for use in electrical circuit breakers typically require two terminal cage screw busbars. Conventional terminals in the electric circuit breakers typically require a terminal cage screw busbars without the possibility to have a quick insertion connection, neither have both types; cage terminal for cables or screwed busbar and quick connection.
One solution to the problem of a two terminal cage screw busbars is disclosed in International Patent Publication No. WO03/028162 (Assigned to Hager Electro S. A.), which uses a fork busbar. According to the Hager patent publication a double connection terminal strip for a modular electrical apparatus is provided with two superimposed levels for connecting conductors capable of being used independently of each other, a first terminal strip of the cage type co-operating with a terminal and a member for clamping/releasing the conductor, accessible from the front surface of the housing of the apparatus including a control lever, the second terminal strip being of the fast-connect type with elastic elements and a without clamping member, the second terminal strip being arranged between the front surface for access to the clamping/releasing member and the cage for connecting the first terminal strip, and being self-supporting by the latter so as not to transmit stresses to the housing.
The present disclosure also provides many additional advantages, which shall become apparent as described below.

Multi-terminal block for electronic devices having superimposed conductor connecting levels


A solid state joint and a method of making a solid state joint between aluminum or magnesium alloys is provided. The joint consists of a diffusion bond with unmelted pieces of a fragmented foil dispersed along the diffusion bond. The joint is made by placing a friable foil between the parts being joined and pressing the parts together to crack the foil. The assembly is then heated and pressed together for a sufficient time, temperature, and pressure to cause the aluminum or magnesium alloy to flow into the cracks, across the foil, and then to diffusion bond together. The foil is made from a material which is harder than the alloy being joined and which will not melt at the bonding temperature.

This invention relates to the field of joining metals by solid state bonding (diffusion bonding) and particularly to joining aluminum alloys, magnesium alloys, or composites having these alloys as matrices.
Solid state bonding, or diffusion bonding, is the metallurgical joining of similar or dissimilar metals by applying heat and pressure for a time sufficient to cause commingling of atoms at the joint interface. The bonding is accomplished without any melting at a temperature which is at or above about one-half the absolute melting temperature of the metal being joined.
Structures have been produced successfully by diffusion bonding titanium alloys using methods as described in U.S. Pat. Nos. 3,920,175 and 3,927,817. Titanium is unique in that it can dissolve its own oxide at the diffusion bonding temperature. Thus, commingling of the metal atoms to form a diffusion bond can proceed without being blocked by the oxide which naturally forms on most metal surfaces. As illustrated in the above mentioned patents, diffusion bonding has been combined with superplastic forming to produce many high quality monolithic aerospace structures.
Unfortunately, the oxide which occurs on aluminum and magnesium alloys is not dissolved into the base metal at diffusion bonding temperatures. These alloys have a high affinity for oxygen and an oxide film is ever present on their surface. If the oxide film is broken under atmospheric conditions, a new oxide layer is formed almost immediately.
Several approaches have been tried to diffusion bond aluminum. In one approach, the oxide is completely removed and the oxide-free surface is kept in a vacuum to prevent reformation of the oxide during bonding. U.S. Pat. No. 4,483,478 describes such a method whereby "sputtering" is used to remove the oxide and a vacuum is maintained during bonding. However, an extremely good vacuum (less than 10.sup.-8 torr) is required to prevent oxide formation on aluminum within a short time. This level of vacuum is not easy to achieve and the need to physically manipulate workpieces inside a vacuum chamber increases the expense and limits the capability of the process.
In another prior art approach, thoroughly cleaned aluminum sheets are compressed (or rolled) together by about 60% strain (i.e., their surface area is increased 60% while in intimate contact). The broken oxide skin rolls into globules generating good metallic contact and a diffusion bond. Although some configurations can be made by this method, the large amount of deformation required is not suitable for many applications because the regions of desired bond location can shift considerably during joining.
Other attempts to bond aluminum use intermediate materials such as coatings or a layer of a bonding material between the parts. Frequently, the intermediate material itself is difficult to bond to the aluminum. If it has a low melting or softening temperature, the resulting bond will be suitable for subsequent superplastic forming which is done at an elevated temperature. Further, in high strength aluminum alloys having a composition close to their solid solubility limits, it is nearly impossible to diffuse the constituents of the intermediate layer back into the alloy itself. These approaches have not provided consistent high strength bonds.

Solid state joint between aluminum alloys and/or magnesium alloys, and a method of making same


When image information is to be transferred between processing units having different image information forms, e.g., the pixel width and plane width of image information, the image information is converted by a reformatting unit using the form attribute information of a transfer source and that of a transfer destination. With this operation, a long processing time is not required to transfer image information between arbitrary processing units having different image information forms, e.g., the pixel width and plane width of image information, and there is no need to impose limitations on the types of forms which can be processed on one apparatus, thus enabling high-speed processing and realizing a flexible system.

1. Field of the Invention
The present invention relates to an image information processing apparatus for transferring, e.g., color image information between different types of processing means and processing the information.
2. Description of the Related Art
Recently, an image information processing apparatus such as an electronic filing apparatus has been put into practical use. In this apparatus, image information such as documents produced in large quantities are read by a two-dimensional scanning unit (scanner) upon optical, two-dimensional scanning, and the read image information is stored in an optical disk. Arbitrary image information stored in the optical disk is retrieved and read out. The readout image information is then output from an output unit, e.g., a CRT display unit or a printer, to be visualized.
Expression schemes for color image information employed by such apparatuses are classified into various types according to storage forms. There are, for example, a plane independent scheme shown in FIG. 1A and a packed pixel scheme shown in FIG. 1B. In the plane independent scheme, R, G, and B components of each pixel constituting image information are respectively stored in storage means 1, 2, and 3 in units of R, G, and B planes as pixel color components, i.e., in units of colors. In the packed pixel scheme, R, G and B planes are sequentially stored in units of pixels. However, in either the plane independent scheme o the packed pixel scheme, image information varies in form, e.g., the data width and number of constituent bits of each pixel and plane to be stored in a storage means.
In the conventional apparatuses, therefore, when image information is to be processed (transferred) between a plurality of processing means, e.g., a scanner, a memory, and a printer, which process image information in different forms, the form of the image information is changed to conform to the form of each means upon arithmetic processing by a general CPU or the like.
For this reason, a long processing time is required to transfer image information between processing means which process image information in different forms, thus inevitably imposing limitations on the types of forms of image information to be processed in one apparatus. Moreover, transfer of image information between arbitrary units which process image information in different forms is inefficient and hence cannot be practically performed.

Color image information processing apparatus capable of high-speed processing image information in different form


A portable balloon tying device having a vertically extending tube with a first end and an opposing second end joined by a transition and an elongated channel extending outwardly as a cantilever from the transition. The first end defines a lower extent of a radiused arcuate surface for being received against a web of a hand between adjacent fingers of a person for using the portable balloon tying device, to cushion a loading and movement of the portable balloon tying device against the web which reduces fatigue and hurting during balloon tying use.

Balloons with ribbons and the like are generally tied with a half-hitch knot. First, a loop is made in the balloon tail, passing the end over the standing part, then passing the end under the standing part and through the loop. See The World Book Dictionary, Clarence L. Barnhart and Robert K. Barnhart, editors, World Book, Inc., publisher, 1990, page 956, column 3. The loop is usually made around one or more fingers, which has disadvantages, among which are: 1. Difficulty of passing the end through the loop, because the material is generally kept very tight to prevent air or gas from escaping from the balloon; 2. Damage to the material while tying the knot therein; 3. Fatigue of the hands, especially the fingers, due to tying many balloons; 4. Extra time is often required to avoid or cope with the above listed disadvantages.
There are various commercially available devices which can be attached to the open ends of balloons to seal the air or gas there inside. Such products do not relate to the invention because when such devices are removed from the balloon, the air or gas escapes therefrom.
Other devices, such as that disclosed by Peverley in U.S. Pat. No. 4,989,906, issued on Feb. 5, 1991, attach to fixed support means via a bracket. Such devices are not generally portable since they must be secured to a fixed support. There is a need for a balloon tying device which can be held in the hand or positioned on a stanchion while operated with both hands to facilitate easy tying of balloons.
U.S. Pat. No. 7,967,344 shows a portable balloon tying device. While this device provided a great leap in the ability to tie a balloon, improvements to this type of device could improve the ability to quickly and efficiently tie balloons so that the balloon did not fall off the end of the device during the tying process. Conversely, an improvement could enable an overly stretched balloon to more easily be removed from the device.
While the foregoing discloses a portable balloon tying device that successfully fixes a knot in a neck portion of an inflated balloon, a need exists to resolve an operational drawback that arises from repeated operation of the device by a balloon inflatest. The drawback is the bearing load and movement of the base against a web portion of the hand during operation of the portable balloon tying device. As explained, the base of the device slides over two fingers of one hand, for example the first and second fingers lateral of the thumb. This disposes an end wall of the base between the second and third fingers lateral of the thumb, with a lower edge bearing on the web of the hand between the adjacent fingers. The loading and movement occurs as a balloon neck is stretched laterally from one side to the other of the open cantilever, wrapped around the cantilever to cross over the stretched neck, and looped over and pulled longitudinally under to form the knot and dislodge the knotted neck from the cantilever. The loading and movement on the web fatigues the hand while holding the device and forming the knot in a plurality of balloons in an extended period, and such may bruise and hurt the hand of the balloon tying operator. This therefore limits the number of balloons that can be comfortably knotted before a rest period may be necessary.
Accordingly, there is need in the art for an improved balloon tying device providing a more comfortable fitting of the base to the hand to reduce fatigue and bruising. It is to such that the present invention is directed.

Portable balloon tying device with reduced fatigue-inducing distal edge


This invention relates to a combustion chamber and a method for burning waste material, such as waste oil, hazardous chemicals, and municipal garbage or trash, in a manner which generates very little smoke, particulate matter, unburned hydrocarbons, or potentially hazardous fumes in the exhaust gases. The combustion chamber is enclosed within an outer wall which provides an inlet for air, fuel, and burnable waste material, and an outlet for exhaust gases. Enclosed within the outer wall is an annular wall, which divides the chamber into a central burning region inside the annular wall, and an annular space between the outer wall and the annular wall. The central burning region contains a primary burning zone, while the annulus space carries a portion of the inlet air along the length of the chamber.

This invention involves a combustion chamber designed to burn waste materials, including hazardous chemicals.
Various types of combustion chambers have been developed to burn waste materials such as waste oil, hazardous chemicals, and garbage. One of the primary reasons for burning waste material is simply to dispose of it; however, the burning of waste may also serve as a source of usable heat (for example, it can be used to heat buildings or to generate steam in a manufacturing or chemical processing facility). As used herein, the term "burning" is interchangeable with "combustion."
It is widely known that complete combustion, which reduces air pollutants and increases energy output, can be promoted by increasing the mixing of the burning material with an oxygen-containing gas such as air or purified oxygen. Along those lines, various devices and components have been incorporated into combustion chambers to increase the air-fuel mixing inside those chambers.
For example, the combustion chamber shown in U.S. Pat. No. 2,072,731 (Crosby, 1937) contains two burning zones, which can be regarded as a primary burning zone and a secondary burning zone. Fresh air is slightly compressed by a fan near one end of the chamber. Part of the air enters a central burning chamber through a first set of inlets. The central burning chamber contains a spray nozzle near the inlet, which injects the fuel into the chamber in small droplets. The fuel is ignited by a spark plug and burns, thereby creating a primary burning zone near the spray nozzle. The remaining portion of the fresh air enters an annulus which surrounds the central chamber. The air passing through the annulus is heated by the high temperatures in the central burning chamber. After being heated, the air enters the central chamber through orifices near the outlet of the device, which pass through the wall that separates the central chamber from the annulus. A secondary burning zone is created when the heated air enters the central chamber and mixes with the partially burned fuel.
Another combustion device is described in U.S. Pat. No. 2,059,523 (Hepburn et al, 1936). In that device, outside air is drawn into a secondary burning zone by means of a draft-producing device at the end of the combustion tube.
U.S. Pat. No. 4,054,028 (Kawaguchi, 1977) describes a combustion device similar to the device of Crosby, except that the air which enters the annulus does not enter the internal burning chamber through numerous small orifices. Instead, it enters the chamber through a small number of "scoops" placed in two distinct stages. Fresh air enters through a primary inlet which surrounds a fuel spray nozzle. The primary burning zone is immediately downstream from the spray nozzle. The secondary burning zone is created by a first set of scoops, which convey heated air from the annulus into the central burning chamber. A second set of scoops also serve as passageways for fresh air; however, those scoops do not create a tertiary burning zone. Instead, they create a dilution zone which assertedly controls the temperatures generated in the combustion chamber, and which presumably reduces the temperatures of the exhaust gases.
The Kawaguchi system can also contain various devices to promote air-fuel mixing, such as (1) angled blades in the primary air inlet, to generate a swirling action in the primary burning zone; (2) baffle disks between the primary and secondary burning zones, to generate turbulence and mixing; and (3) air scoops which are slanted in the upstream direction to generate turbulence and mixing in the secondary burning zone.
U.S. Pat. No. 4,173,118 (also by Kawaguchi, 1979) discloses a similar device with (1) numerous fuel spray nozzles and air inlets, spaced in a circle to create a primary burning zone with an annular rather than circular shape; (2) a water or steam injector in the middle of the primary burning zone, to aid the burning of heavy oil. That system also contains a secondary burning zone, and a dilution and cooling zone, similar to U.S. Pat. No. 4,054,028.
U.S. Pat. No. 2,398,654 (Lubbock et al, 1946) and U.S. Pat. No. 4,819,438 (Schultz 1989) also describe a combustion chambers with primary and secondary burning zones followed by a dilution zone. In the Lubbock patent, the air in the annulus flows in a direction counter to the flow of the exhaust gas, which causes it to be heated before it enters the primary burning zone. In the Schultz patent, steam is circulated around the burning chamber to control the temperature and to obtain a useful work output.
U.S. Pat. No. 3,567,399 (Altmann et al, 1971) discloses an afterburner. Hot waste gases from a primary burning chamber are introduced along with an additional source of fuel such as natural gas which is ignited near the inlet. Air is compressed into an annulus surrounding the afterburner chamber, and enters the afterburner through various orifices which are preferably angled to promote a swirling motion inside the afterburner.
The combustion chambers described above which are designed for burning waste material suffer from various limitations. The more complex devices, such as the burning chambers invented by Kawaguchi, must be assembled from dozens of component parts, which must be manufactured with close tolerances at relatively high expense. In addition, a complex burning devices can be difficult to maintain properly, especially if used to burn waste material having a wide variety of characteristics. The simpler devices such as the burners invented by Crosby, Lubbock, and Schultz do not ensure the desired degree of combustion in some situations, and can cause the emission of smoke, unburned hydrocarbons, and other air pollutants. In addition, none of the devices above are suitable for burning solid wastes.
Other devices have been developed to promote air-fuel mixing in gas turbines and jet engines; for examples, see U.S. Pat. No. 3,593,518 (Gerrard, 1971), U.S. Pat. No. 3,937,007 (Kappler, 1976), and U.S. Pat. No. 2,930,194 (Perkins, 1960). Those are not relevant to the subject invention, for several reasons. Gas turbines and jet engines involve extremely high gas velocities, the internal mechanisms involve rapidly-moving and often delicate components, they are usually operated at very high levels of excess oxygen, and they cost a great deal more to manufacture than waste oil burners.
There is, therefore, a need for a sturdy and inexpensive device for burning waste material, which has no internal moving parts and which can achieve higher levels of combustion of burnable material than previously achieved by other combustion chambers.
One object of this invention is to provide a simple, sturdy, inexpensive, and reliable device for burning waste material.
Another object of this invention is to provide a waste-burning device which can achieve virtually complete combustion of any liquid or liquified input material, such as waste oil, paint sludge, or hazardous chemicals.
Another object of this invention is to provide a device for achieving virtually complete combustion of hazardous chemical wastes, to ensure that very low quantities of toxic, carcinogenic, or otherwise hazardous air pollutants will be emitted into the atmosphere.
Another object of this invention is to provide a device for achieving virtually complete combustion of solid wastes such as old tires and municipal or industrial trash or garbage.

Waste combustion chamber with tertiary burning zone


A computer network provides distribution of a multithread process among different data processing stations in the network. In one embodiment, a global name server is used to distribute the threads and update values associated with the threads. When the process is created, the global name server establishes a distributed process context by dividing logical addresses used by the process among physical addresses corresponding to respective portions of the system memory devices (RAM) in different data processing stations. The global name server assigns a universally unique identification number (UUID) to each process, and uses the UUID to manage the distributed threads. When the data in the distributed process context is updated, the changes are apparent to all threads executing in the distributed process. This approach simplifies access and retrieval of shared process context information, and may be implemented using current thread models, making it easier to incorporate into existing networks.

1. Field of the Invention
The present invention generally relates to computer systems, and more particularly to a method of distributing a multithread process among multiple workstations of a computer network, for improved high performance multithread data processing.
2. Description of the Related Art
The basic structure of a conventional computer system includes one or more processors which are connected to several input/output (I/O) devices for the user interface (such as a display monitor, keyboard and mouse), a permanent memory device for storing the computer's operating system and user programs (such as a magnetic hard disk), and a temporary memory device that is used by the processors to carry out program instructions (such as random access memory or RAM). The processors communicate with the other devices by various means, including a system bus or a direct channel.
When a user program runs on a computer, the computer's operating system (OS) first loads the main program file into system memory. The program file includes several objects (values) stored as data or text, and instructions for handling the data and other parameters which may be input during program execution. The processors use "logical addresses" to access the file objects, and these logical addresses correspond to physical addresses in RAM. Binding of instructions and data to physical memory addresses is accomplished by compiling the program file using relocatable code, which is indexed (linked) to physical memory by the OS loader during loading of the file.
A computer program can be broken down into a collection of processes which are executed by the processor(s). A process is a set of resources, including (but not limited to) logical addresses, process limits, permissions and registers, and at least one execution stream. The smallest unit of operation to be performed within a process is referred to as a thread. The use of threads in modern operating systems is well known. Threads allow multiple execution paths within a single address space (the process context) to run concurrently on a processor. This "multithreading" increases throughput in a multiprocessor system and provides modularity in a uniprocessor system.
In a single tasking operating system, a computer processor executes computer programs or program subroutines serially, that is no computer program or program subroutine can begin to execute until the previous computer program or program subroutine has terminated. This type of operating system does not make optimum use of the computer processor in a case where an executing computer program or subroutine must await the occurrence of an external event (such as the availability of data or a resource) because processor time is wasted. This problem lead to multitasking operating systems wherein each of the program threads performs a specific task. If a thread being executed must wait for the occurrence of an external event, i.e., the thread becomes "non-dispatchable," then its execution is suspended and the computer processor executes another thread of the same or different computer program to optimize utilization of processor assets. Multitasking operating systems have also been extended to multiprocessor environments where threads of the same or different programs can execute in parallel on different computer processors.
FIG.1 illustrates multithreading in a prior art multiprocessor computer system 10. System unit 10 includes one or more processors 12a, 12b and 12c which are connected to various I/O devices 14 and a memory device 16 (RAM) via a bus 18. Each processor includes a central processing unit (CPU) and one or more on-board caches. Each cache comprises a small amount of high speed memory which stores a local copy of data utilized by its associated processor. A typical processor used for multithread processing is the PowerPC.TM. integrated circuit superscalar microprocessor manufactured by International Business Machines Corporation. When data requested by a processor is not resident within its associated cache, the processor will attempt to load the requested data from an optional (L2) cache, or from global memory 16 (which may include one or more individual modules of physical memory).
Global memory 16 has a kernel portion with a set 20 of thread context fields for N number of threads associated with a particular process. Global memory 16 further includes a process context 22 in a user address space which contains all of the logical addresses for data and instructions used by the process. After a thread is created and prior to termination, the thread will most likely utilize system resources to gain access to process context 22. Through the process context, process threads can share data and communicate with one another in a simple and straightforward manner.
Separate computers can be linked into vast networks which allow for the distribution of processing among various systems in the networks. Unfortunately conventional systems only support distributing entire processes across the network, and do not support distribution of threads, since there is no common process context for the various systems in the network. This limitation causes programmers to use contorted mechanisms when writing distributed programs that need to work together to accomplish a task. These mechanisms often involve complicated communication protocols (such as sockets or pipes) that are very time consuming to develop and difficult to debug. An alternative approach is to provide a common process context using a shared virtual memory system which maps a large virtual memory space globally to all processes in the system. This approach, however, causes problems in performance, security and integrity. It would, therefore, be desirable and advantageous to provide an improved method to allow execution of various threads of a single process on different systems in the network.

System for distributing a plurality of threads associated with a process initiating by one data processing station among data processing stations


A system includes at least one application device, a control unit, at least one detection device, and at least one processor configured to be disposed on-board a vehicle system. The at least one application device is configured to be at least one of conductively or inductively coupled with at least one conductive track of a route traveled by the vehicle system. The control unit is configured to control the at least one application device to electrically inject at least one examination signal into the conductive tracks. The at least one detection device is configured to detect the examination signal passing through a test loop. The at least one processor is configured to identify a failure of the vehicle to adequately shunt electrical signals between the conductive tracks based upon the examination signal detected by the at least one detection device.

Various transportation systems, such as signaling systems, or systems that operate highway crossings, may detect the presence or approach of a vehicle (e.g., rail vehicle) to operate signals and/or crossings appropriately, and/or to schedule vehicle travel through a signal and/or crossing. Rail vehicles may be detected using a principle that axles of a rail vehicle that are electrically conductive will effectively short two rails of a track together. This shorting of two rails via an axle of a vehicle traveling over the rails is referred to herein as shunting. This shunting may de-energize track circuits for operating signals, or may be monitored for operating highway crossings. A track signal may be transmitted, for example, proximate to a signal, and that track signal may be received at a distance from the signal. A track signal may be transmitted, for example, proximate to a crossing, and a permanent shunt across the rails may be disposed at a distance from the crossing corresponding to an expected or permitted vehicle speed. The higher the speed, the farther away the permanent shunt may be. A track signal may be sent through a rail, with the signal traveling through the rail, through the track receiver or permanent shunt, and back toward the source of the signal. When a rail vehicle approaches the signal or crossing and crosses the permanent shunt so that the vehicle (or a portion thereof) is closer to the signal or crossing than the track receiver or permanent shunt is, the track signal is then routed back to the source via an axle (or axles) of the vehicle instead of via the permanent shunt. By monitoring, for example, the signal at the track receiver, the location of the vehicle between the transmitter and receiver may be detected, and a signaling system may be appropriately controlled. By monitoring, for example, a rate of change of electrical impedance of the circuit between the transmitter and the shunting, the approach of the vehicle within the distance to the permanent shunt may be detected, and a highway crossing system may be appropriately controlled.
However, the axles of a rail vehicle may not provide the desired shunting used to determine if the rail vehicle is within a predetermined range, for example, of a crossing or signal. For example, a running surface of the route (e.g., rails) and/or treads of wheels may be contaminated or covered by a non-conducting substance, such as a rust film. Shunting may be prevented until the non-conducting substance wears off, or wears off enough to allow shunting. If the rail vehicle fails to shunt, the presence of the rail vehicle may not be detected by crossing or signaling systems, resulting in delay, inconvenience, and/or reduced safety. Conventional approaches may use human inspectors who move along the track, for example, infrequently used lines, to inspect for rust or other material that may inhibit shunting. This manual inspection is slow and prone to errors.

System and method for shunting detection