A fiber optic refractometer is disclosed for use in on-line measurement of the refractive index of a process fluid. The refractometer does not require light to pass through the process fluid and is hence unaffected by the presence of light diffusing particulate matter in the process fluid. A particular application for which the refractometer is well-suited is the on-line measurement of the hydrogenation state of edible oils (which correlates with refractive index) during the partial hydrogenation process.

Edible oils are those extracted from plant seeds, such as soybean, cottonseed, peanut, sunflower, corn, etc., and processed for human consumption. They are usually found as ingredients in foodstuffs such as bread, pastries, margarines, salad dressings and candies. Today, they are consistently used as the medium to fry foods. In practically all cases, these oils undergo the process of partial catalytic hydrogenation to prepare them for use in their intended product. Hydrogenation alters the molecular structure within the fatty ester chains of the oil triglyceride molecule by reducing the number of double carbon bonds that naturally occur.
Partial hydrogenation affects edible oils in two ways. First, it stabilizes the oil; it extends the time that the oil's flavor and smell are acceptable as a human food. Second, it converts the naturally occurring oils into forms that have melting and handling characteristics demanded by the application; the more an edible oil is hydrogenated, the higher its melting point becomes. For example, oils used in salad dressings must not coagulate on cold lettuce; margarine should remain firm on the breakfast table, yet melt on pancakes and biscuits; shortening should remain firm in the pie crust or cake batter, but quickly melt in the oven.
The basic batch hydrogenation process in use today has changed little since it was invented nearly a century ago. A closed vessel (converter) is filled with approximately 58,000 pounds of oil. The oil is heated to approximately 175.degree. C. with steam passing through internal coils. Then roughly 10.5 pounds of nickel powder is added as a catalyst; hydrogen gas is introduced continuously into the bottom of the converter and the hydrogen is vigorously mixed into the oil and nickel powder by an agitator. The ensuing hydrogenation reaction is exothermic and is controlled by cooling water passing through the original steam coils. The converter operates at a pressure of 20 to 30 psi while still releasing excess hydrogen and processing effluents from the top of the converter. Hydrogenation is arrested immediately upon cessation of agitation which is normally done several times during the process. Processing times vary from 15 to 60 minutes depending on the desired degree of hydrogenation.
Although it is not a part of the actual hydrogenation process, about 30 pounds of diatomaceous earth are also added into the converter vessel along with the nickel catalyst. The diatomaceous earth is used as a filtering aid to remove the nickel from the oil at the termination of the process. The finished batch of oil is placed in a holding tank and continuously circulated through wire filter screens to cleanse the oil of the nickel catalyst. The diatomaceous earth is trapped on the wire screen filters, in effect becoming the filter itself. The nickel powder is in turn filtered out of the oil by this diatomaceous earth "cake." The "cake" is reclaimed off of the filter screens and the nickel recycled. All in all during the actual hydrogenation process, there is a 0.072% mixture (by weight) of nickel catalyst and diatomaceous earth present in the oil.
There are several ways of accurately determining the amount of hydrogenation of edible oils. One is the Wijs method, recommended by the American Oil Chemists, Society (AOCS). A large, measured amount of iodine monochloride reagent is added to a specific quantity of partially hydrogenated oil. The reagent "saturates" the oil by eliminating any double carbon bonds remaining in the various esters. The excess reagent is titrated with sodium thiosulfate, using starch as an indicator. The amount of iodine reagent lost is called the iodine value, a number that decreases with increased hydrogenation. This method, though accurate, is a laboratory procedure requiring 30 to 45 minutes to perform.
Another technique of determining the degree of partial hydrogenation is to measure the oil's refractive index (RI). There is a direct, linear relation between RI and iodine value. RI linearly decreases with increased hydrogenation. An example is the relationship between RI and iodine value (I) of cottonseed, peanut, soybean, and linseed oils. It has the form EQU n.sub.D.sup.40 =1.4515+(0.000117)I (1)
where n.sub.D.sup.40 is the RI at a wavelength of 589.3 nm and at 40.degree. C. (See Bailey, A. E., Bailey's Industrial Oil and Fat Products, John Wiley & Sons, Inc., Vol. 2, 1982, pp. 40, 220). The AOCS Method Cc 7-25, "Refractive Index," is an accepted industry standard in determining the degree of edible oil hydrogenation. The technique is quick and can be performed in the process control room.
Several established designs of process refractometers are in common use today which depend upon light transmission through the oil in order to measure its RI. With any of these units, a reading can be made in a matter of minutes and a correlation of RI to iodine value read from a chart. The handicap with this method is that these refractometers operate on light transmission through the oil sample, and with nickel catalyst and diatomaceous earth in the oil, the transmitted light is diffused, making the refractometer response difficult to interpret. Therefore, an oil sample withdrawn from a batch of oil being processed must first be manually filtered before analysis. While the sample is being withdrawn, filtered and analyzed, the hydrogenation process must come to a standstill which sometimes means a delay of 10 to 20 minutes. It is common to have two or three test samples taken during the processing of an oil batch. Not only is this a worrisome delay to the oil processor, but also the accuracy and repeatability of the refractometer reading is subject to operator skill. It would be desirable to have an on-line refractometer that will perform accurately in the presence of catalyst and diatomaceous earth.
One type of refractometer that does not require light transmission through the process liquid, and hence is unaffected by the presence of catalyst and diatomaceous earth, is a critical angle refractometer. (See Maley, L. E., "Refractometers," Journal of Chemical Education, vol. 45, no. 6, 1968, pp. A467-A485.) Such a device consists of a glass prism or plate with a flat side mounted over a small opening in a pipe or tank wall. The back side of the glass (RI=n.sub.g) is in direct contact with the process liquid (RI=n.sub.1). Using Snell's law, a critical angle, .theta..sub.c, for light incident at the glass-liquid interface is selected according to ##EQU1## such that light is reflected back to a sensor when EQU n.sub.1 .ltoreq.n.sub.g sin .theta..sub.c ( 3)
While the operation of this device is not affected by the turbidity of the process liquid, the device is only a fixed comparator and is thus unsuitable for continuous measurement of refractive index values. Also, n.sub.g must be greater than n.sub.1 for there to be a critical angle.
Another type of refractometer which does not require light to pass through the process liquid, but offers a means of continuously measuring RI values, is a fiber optic refractometer. The efficiency with which optical fibers transmit light is determined by the disparity of RI that exists between the core and cladding materials. It follows that such a device could be used as a refractometer if the process liquid of interest became the "cladding" about a glass core. By measuring the efficiency at which such a fiber transmitted light energy, the RI of the liquid cladding could be determined. The concept of attenuated total reflectance (ATR) forms this major category of fiber optic refractometers. (See Kapany, N. S., and J. N. Pike, "Fiber Optics, part IV, a photorefractometer," Journal of the Optical Society of America, vol. 47, no. 12, 1957, pp. 1109-1117.) In these instruments, a conical beam of light with a uniform intensity, I watts/steradian, excites a glass rod or fiber (RI=n.sub.g). The transmitted light is then measured by a photosensitive device. The following equations define the numerical aperture NA, and transmitted light power, Pt ##EQU2## By measuring Pt, continuous values of n.sub.1 may be calculated provided n.sub.g and I are known.
In order for the device to work, n.sub.1 must be less than n.sub.g. To overcome the problem of the refractometer losing all resolution when n.sub.1 .gtoreq.n.sub.g, glass rods or fibers of a higher RI must be substituted.
A variation of the ATR fiber optic refractometer uses a laser beam incident on the end of the glass rod. (See David, et al., "Design, development and performance of a fiber optics refractometer: Application to HPLC," Review of Scientific Instruments, vol. 47, no. 9, 1976, pp. 989-997; also, U.S. Pat. No. 3,999,857, J. D. David, D. A. Shaw & H. C. Tucker, "Refractive Index Detector," 1976.) The beam angle into the rod is adjusted via a mirror moved by a micrometer until the edge of the "cone of acceptance" (i.e., the numerical aperture or NA) is found. Multiple reflections of the light propagating down the fiber make the transition very sharp. The micrometer reading correlates to the NA, and n.sub.1 can be calculated from equation (4a). The instrument locates the sharp light transition at the edge of the NA, but its output drops to a low, constant level once the incident beam angle exceeds the NA.
A fiber optic refractometer using Fresnel's equations has also been designed. (See Meyer, M. S., and G. L. Eesley, "Optical Fiber Refractometer," Review of Scientific Instruments, vol. 58, no. 11, 1987, pp. 2047-2048.) Monochromatic light is transmitted down a single mode fiber and reflects off the far end of the fiber, immersed in the process liquid. The core (RI=n.sub.co) at that end of the fiber is polished smooth, perpendicular to the fiber axis. Fresnel reflections from the core/liquid dielectric interface are transmitted back through the fiber to a photosensor. For light perpendicularly incident to such an interface, the Fresnel's reflection is ##EQU3## The reflected light is correlated to equation (5), and the refractometer operates for values of n.sub.1 either greater than or less than n.sub.co. There is a reflection null at n.sub.co =n.sub.1. However, due to the quadratic dependence of the reflectance, it is impossible to tell which side of the reflection null the refractometer is operating on. The solution is to use two fibers of different RI and observe their relative outputs.
Fiber optic refractometers using bent fibers have also been developed. (See Golunski, W., et al., "Optical fiber refractometer for liquid refractive index measurement," Proceedings of the SPIE--Optical Fibers and Their Applications V, vol. 1085, 1990, pp. 473-475.) Bending the fiber effectively decreases n.sub.co, reducing the NA of equation (4a). Consequently, the fiber's transmitted power, Pt, of equation (4b) is also reduced. One paper, using a stripped, step index quartz fiber, describes how the range of such a refractometer's operation (analog range of Pt) can be varied by controlling the radius of the fiber bend; the smaller the radius, the lower the RI sensing range becomes. (See Harmer, et al., "Optical fibre refractometer using attenuation of cladding modes," Proceedings of First International Conference on Optical Fibre Sensors, Electronics Div., Institute of Electrical Engineers, 1983, pp. 104-108.) Such a refractometer is only functional when the liquid RI is less than the effective n.sub.co.
Previous fiber optic refractometers, as described above, suffer from several shortcomings. First, analog optical responses are subject to variations in light source intensity and uniformity as well as light sensor sensitivity, liquid turbidity, operating temperatures and the bending of the optical fiber. Second, most responses of Pt versus liquid RI are nonlinear and must be compared to a complicated algorithm in order to maintain high accuracy. Third, there is no compensation for RI dispersion due to the light source wavelength. Fourth, mechanical measurement devices are affected by misalignment.


A compensating accelerometer includes a housing, and a sensing element mounted in the housing. The sensing element includes a pendulum flexibly mounted on a base. The sensing element includes a coil mounted on a movable plate, a curtain having a slit, and a load mass affixed to the curtain. An angle sensor within the housing includes the curtain, a fork, a light source and a differential light detector. A momentum sensor within the housing includes a permanent magnet, an inner magnetic conductor, an outer magnetic conductor and the coil. The permanent magnet is magnetized in a direction of an axis of the pendulum. A stopper is used to adjust position of the fork. A spring is on opposite side of the fork, for taking up slack in the fork. A fixator has an eccentric, for adjusting a gap between the curtain and the light detector.

1. Field of the Invention
The present invention relates to measurement equipment, and more particularly to compensating pendulum-type linear accelerometers with flexible suspension of the sensing element.
2. Description of the Related Art
An example of a conventional accelerometer is illustrated in Japanese Patent No. 188924/1981. This accelerometer includes an inertial mass (in this case, a pendulum) that is placed in a sealed housing. Inside the sealed housing there is typically a vacuum or an inert gas, such as helium. The inertial mass is generally in the shape of a movable rod, whose lower portion is attached, using flexible suspension, to a plate. The plate is then affixed, using screws, to a base surface of one of the frames. The rod can rotate, using the flexible suspension, around the plate in the direction of the measurement axis of the accelerometer. The rod of the pendulum has two cylindrically shaped coils attached to it.
The axial magnetic systems of the accelerometer include magnetic conductors, in this case first and second frames, which are made of magnetically soft material. The magnetic axial systems also include permanent magnets and field concentrators, such that the movable coils are located in the gaps formed between the surfaces of the magnetic conductors and the field concentrators. The above elements together with the movable coils comprise the momentum sensor of the accelerometer.
Both frames are rigidly connected to each other, and each of the frames includes stoppers for limiting the range of motion of the pendulum. This range of motion can be regulated by moving the stoppers in the threaded openings in the frames.
Inside the housing there is an angle sensor that measures the movement of the pendulum. The angle sensor consists of a single light source and two light detectors. The pendulum is located between the frames, such that the light source is located on one side of the pendulum, and the light detectors are located on the other side of the pendulum. Thus, the pendulum acts as a shading element that blocks the light from the light source to the light detectors. When there is no measured acceleration along the measurement axis, the pendulum is in a neutral position, and some of the light from the light source is shielded by the pendulum, while the remainder of the light is evenly distributed between the two light detectors. Thus, the current produced by the light detectors (if, for example, photodiodes are used) is equal. Note that the surfaces of the photodiodes (light detectors) need not be entirely illuminated, but so long as the illumination of each detector is the same, the output currents are the same, indicating that the measured acceleration is zero.
When the sensor experiences acceleration along the measurement axis, the pendulum is displaced from its neutral position due to inertia force. As a result, the light distribution between the two detectors changes, and therefore the relationship between the current from the two detectors also changes. The difference between the two currents is related to the acceleration, and using conventional electronics, can be converted to an acceleration value. The current can also be used in a feedback circuit; it passes through the movable coils interacts with the magnetic field of the permanent magnets and returns the pendulum to its neutral position. The magnitude and polarity of the feedback current therefore permit a measurement of the acceleration.
The conventional accelerometer described above has a number of disadvantages. One of the disadvantages is that due to the relatively large dimensions of the rod of the pendulum, the light detectors need to be placed relatively far apart, which tends to increase noise and reduces the sensitivity of the measurement. Another disadvantage is that the neutral position of the pendulum typically does not precisely correspond to a zero output signal of the accelerometer. Yet another disadvantage is that mechanical tuning of the accelerometer tends to be difficult, particularly with regards to the zero bias signal.
Another example of a conventional accelerometer is described in U.S. Pat. No. 4,649,748. The accelerometer described in this patent shares some commonalities with the earlier-described accelerometer. However, the accelerometer in U.S. Pat. No. 4,649,748 has the following differences compared to the device of Japanese Patent No. 188924/1981: first, the first frame—the magnetic conductor of the magnetic system—is used as a mounting base of the accelerometer, which permits simplifying the manufacturing of the accelerometer, and also reduces its mass.
Additionally, the free end of the pendulum rod, which is located between the light source and the light detectors, is formed as a thin plate which permits the light source and the light detectors to be closer to each other. This increases the sensitivity of the angle sensor, and reduces the noise in its output. One of the options is for the thin plate to also have the shape of a rod, another option is to have the plate shaped as a thin rod with a slit along the axis of the rod, such that part of the light passes through the slit.
The various elements of the overall construction where the light source and light detectors are housed are pressed against the frames using a flat spring, which increases the precision of the mechanical tuning of the zero bias signal.
The device described in U.S. Pat. No. 4,649,748 has a number of disadvantages. For example, one of the disadvantages is due to the use of two magnetic systems. In case of a need to increase range of measured acceleration, this leads to a relatively values for currents that might flow through the coils, which in turn leads to a large amount of heat being dissipated, which in turn leads to an error source in the measurement, relating to the waste heat. Additionally, the use of two different magnetic systems increases the manufacturing cost and complexity of the device. These disadvantages are also present in the device described in Japanese Patent No. 188924/1981, discussed above.
Furthermore, the mechanical tuning of the zero bias signal is relatively coarse. Furthermore, a relatively small range of displacement of the light source and detectors relative to the housing makes it difficult to tune the angle sensor when there is a relatively large angle between the axis of the pendulum rod in the neutral position, and the base mounting surface of the accelerometer.

Compensating accelerometer with optical angle sensing


A horseshoe adapted to be compliant with the expansion and relaxation of the heel areas of a hoof of a horse in motion. The horseshoe consists of a toe area hingedly connected to two heel areas. In the preferred embodiment, the hinges are formed of a resilient material and are tapered from a wide end on the inside of the horseshoe to a narrow end on the outside of the horseshoe. Metal toe and heel inserts are provided as a wear surface for the horseshoe and the toe insert is sufficiently malleable to allow shaping of the horseshoe using typical blacksmithing equipment and tools.

This invention relates generally to horseshoes and more particularly to a horseshoe which is compliant with the relative movements of the hoof of a horse.
Horseshoes have been used for centuries in order to protect the horse""s foot and to enhance performance. Most horseshoes in use today are made of metal such as steel, aluminum alloys, and rarely, stainless steel or titanium. Aluminum alloys are most commonly used today in racing and are characterized by relatively low weight and expense. Horses engaged in training or racing are commonly reshod every four to five weeks. This is done in consideration for the wear incurred by the horseshoes, but also by the desire to maintain healthy geometry with respect to the configuration of the foot and so facilitate optimal biomechanics, or way of going.
The surfaces on which horses train and perform also vary widely. Horses frequently train and compete on grass, sand, cinder, crushed stone, and sometimes on packed surfaces which nearly approach the hardness of asphalt or cement. The hardness of the training or racing surfaces can greatly increase the effective rate of loading, thus the shock and vibration, e.g., the peak g forces which the horse will experience. Such will necessarily influence the nature of the waveform and the primary frequencies of shock and vibration transmitted to a horse""s anatomy. These factors can directly affect a horse""s efficiency, athletic performance and the amount of trauma that will be experienced.
While the characteristic biomechanics or way of going of horses can vary, e.g., as between Standardbreds and Thoroughbreds, or due to different gaits, such as running, jumping, walking, etc., what normally happens as a horse""s foot and hoof impact the ground is that the back of the foot touches first, then the foot flattens and slides anteriorly skating across the surface. In particular, this is true of a horse""s rear foot. Then the rear portion of the foot is loaded and deflection takes place. This will generally cause the foot to rotate backwards at the heel. The foot will then rotate forwards and recover to a relatively neutral position. Subsequently, the foot rapidly rotates forwardly and about a rocker point located between the geometric center of the foot and a short distance behind the anteriormnost area of the toe as the foot breaks over and toe-off takes place, thus ending the ground support phase and beginning the flight phase.
The heel of a horse""s foot normally strikes the ground slightly before the toe, and this results in immediate heel expansion due to the action of the frog. As the frog is forced upward, the frog stay acts as a wedge in the digital cushion. This forces the digital cushion to expand, primarily in the outward direction because it is confined by structures of the foot in the dorsal, volar, and proximal directions. After the hoof is lifted from the running surface, the heel areas contract.
Furthermore, it is known that in the unshod natural state, a horse""s foot and hoof will flex and slightly widen when it is loaded. The use of relatively rigid metal or aluminum horseshoe substantially prevents this natural movement and so tends to reduce both the effective size, and the shock and vibration absorbing capability of a horse""s foot. A steel horseshoe is known to be more flexible in this regard than an aluminum or titanium horseshoe. It is believed that the occurrence of hoof cracks is sometimes caused by the flexing and widening action of the foot and hoof working against the nails associated with a substantially inflexible horseshoe.
One of the challenges encountered when attempting to reduce the rate of loading, and attenuate the shock and vibration experienced by a horse is posed by the fact that a horse is a rather large animal, e.g., commonly weighing between 800-1400 pounds, and when running at speeds between 30-40 miles per hour, a load exceeding 15,000 pounds can be placed upon a horse""s leg. Accordingly, approximately 2600 pounds per square inch can be placed upon a typical horseshoe having roughly 6.5 square inches of working surface. When running on a hard race track, the entire duration of the impact event can be as short as 1.5 milliseconds, and over 350 peak g""s can then be experienced.
The typical metal horseshoes restrict the natural expansion and contraction of the hoof of a horse in motion and can cause increased stress in the hoofs and legs of a horse in motion.
The foregoing illustrates limitations known to exist in present horseshoes. Thus, it is apparent that it would be advantageous to provide an alternative directed to overcoming one or more of the limitations set forth above. Accordingly, a suitable alternative is provided including features more fully disclosed hereinafter.
In one aspect of the present invention, this is accomplished by providing a horseshoe having at least two hinges, a toe area, two heel areas and two quarter areas. The at least two hinges are positioned to allow expansion and relaxation of the horseshoe in compliance with the hoof of a horse in motion. Preferably, the hinges are positioned in the quarter areas of the horseshoe and have a tapered shape with the wide end of the taper being towards the interior of the horseshoe.
The foregoing and other aspects will become apparent from the following detailed description of the invention when considered in conjunction with the accompanying drawing figures.


A method of employing a pressure generating system that provides a pressure therapy to the pulmonary system of a patient and to enable the patient to interact with the pressure support system, for example, to monitor the effectiveness of the treatment provided by the pressure support system. The interactive system includes an output device that provides first information to the patient, an input device that allows the patient to provide second information that is based on the first information and a control unit that controls the operation of the output device to present the first information and collection of the second information via the input device. In a further embodiment, the operation of the pressure generating system is controlled based on the second information provided by the patient.

1. Field of the Invention
The present invention pertains to a method of treating a breathing disorder using an interactive pressure support system, and, in particular, to method of providing a pressure therapy to a patient while also providing an interactive capability, so that the patient can answer questions or perform tasks that ascertain the effectiveness of the pressure therapy in treating the patient's breathing disorder.
2. Description of the Related Art
There are many medical ailments that are diagnosed and/or monitored using a questionnaire, survey, or cognitive test taken by the patient. For example, a common method to determine whether a patient suffers from a sleep disorder, such as obstructive sleep apnea (OSA), is to measure the patient's sleep propensity. The patient's sleep propensity and/or changes in the sleep propensity can also be used to determine the severity of the disorder and/or monitor the changes in the patient's condition. Conventional techniques for measuring a patient's sleep propensity include the use of the Epworth Sleepiness Score (ESS), the Functional Outcomes of Sleep Questionnaire (FOSQ) and the Berlin questionnaire. The ESS is determined based on the patient's retrospective reports of dozing behavior in a variety of situations commonly encountered in normal daily life. These retrospective reports are elucidated from the patient via a series of questions. The patient's responses to these questions are tabulated and used to determine the ESS to evaluate his or her sleep propensity.
Currently, the Epworth Test and FOSQ are typically administered on paper. To do so, the written test must be physically supplied to the patient and collected after the patient completes the questions. The test administrator manually tabulates (or uses a computer to tabulate) the responses provided by the patient and calculates the ESS or other score based on the patient's responses. It can be appreciated that the administrative requirements, such as the distribution, collection, time stamping, tabulation, scoring, storing and record keeping, required by this conventional testing technique place a significant burden on the test givers. This burden increases with the number of patient's taking the test as well as the number of times the test is administered to each patient. Typically the same patient will take the Epworth test multiple times during his or her treatment in order to monitor the effectiveness of the treatment therapy. It can thus be appreciated that patient follow-up to determine, for example, the effectiveness of a therapy intended to treat a sleep or breathing disorder, is a relatively expensive and burdensome process.
Another conventional technique that measures a patient's reaction time and cognitive alertness, which are generally understood to be indicative of a patient's sleep propensity, is the Vigilance test. This test is typically administered using a personal computer (PC) and requires that the patient provide responses via the input devices associated with the PC to displayed indica. The patient's reaction time in providing the response and/or the accuracy of the response to the displayed indica are measured and stored to determine the patient's level of alertness. For example, the patient is shown a recognizable object on the PC display, and the patient's reaction time in identifying the object and the accuracy of the identification are measured.
Because this test typically requires a PC to administer, it can only be administered at a facility where an appropriately programmed PC is located, which requires that the patient travel to a place where such a PC is located. Alternatively, an appropriately programmed PC can be provided to the patient. However, this latter alternative represents a significant cost in furnishing the PC to the patient and requires training the patient on the use of such a relatively complicated device. It can thus be appreciated that either alternative for administering the Vigilance test represents a significant burden on either the patient or the test administrator.

This patent is subject to a terminal disclaimer.

Method of providing interactive pressure support


An SCPA includes a pad, capacitative elements, amplifiers on an IC chip. The capacitative elements are disposed on a first circle whose center is located on the pad. The amplifiers which correspond to the capacitative elements are disposed on a second circle which is a concentric circle larger than the first circle. The pad, each of the capacitative elements, and a corresponding one of the amplifiers are aligned in a line so that the length of wiring is the shortest.

1. Technical Field
The present disclosure relates to a power amplifier and a transmission device which use a complementary metal-oxide semiconductor integrated circuit (CMOS-IC).
2. Description of the Related Art
An exemplary known power amplifier using a CMOS-IC is a switched capacitor power amplifier (SCPA) which achieves high efficiency (for example, see Sang-Min Yoo, Jeffrey S. Walling, Eum Chan Woo, Benjamin Jann, and David J. Allstot, "A Switched-Capacitor RF Power Amplifier", IEEE JOURNAL OF SOLID-STATE CIRCUITS, VOL. 46, No. 12, DECEMBER 2011). An SCPA which is a kind of class D amplifier performs amplitude modulation by changing the number of amplifiers switching between the power supply voltage and the ground (GND).
However, an SCPA of the related art fails to maximize the output power and the power efficiency. This problem will be described by using FIG. 1. FIG. 1 is a diagram illustrating an exemplary configuration of an SCPA of the related art.
In FIG. 1, the SCPA includes an output terminal 10, multiple capacitative elements (C) 11 to 15, and multiple amplifiers (AMP) 21 to 25. The capacitative elements 11 to 15 and the amplifiers 21 to 25 are disposed in such a manner as to be aligned in lines. In this configuration, the lengths of wiring from the amplifiers 21 to 25 to the output terminal 10 are different from each other. For example, the length of wiring from the amplifier 21 to the output terminal 10 is larger than that from the amplifier 23 to the output terminal 10. Therefore, the former wiring has larger parasitic capacitance and larger parasitic resistance. Accordingly, the output power and the power efficiency fail to be maximized.


A ham product which is formed by cutting a spirally sliced half ham longitudinally through the bone. The cut can be made by feeding the half ham through a band saw blade by hand, by pushing a split cart through the blade, or by feeding the product through the blade on a split belt conveyor. A special tray for packaging of the ham product has either a flap providing a double layer of material at the location of a sharp exposed corner of the femur or a plurality of cushioning dimples extruded at appropriate location thereon to pad the femur corner.

Spirally sliced hams have proven to be popular with consumers, largely because of the convenience of being able to easily remove the presliced ham pieces from the product at the point of use. The ham is produced by slicing it in a manner to form a continuous spiral cut around the center, often with the femur remaining in place but sometimes without a bone. In a "bone in" product, the meat remains continuously attached to the bone all the way around in a spiral pattern adjacent to the spiral cut. Thus, the meat remains connected to the bone firmly enough that the product can be handled and packaged intact without pieces falling off, folding over, or otherwise deforming. At the same time, the slices can be easily detached from the bone at the point of use.
It has been recognized that many consumers are hesitant to purchase an entire ham because there is more meat than the household of the purchaser can conveniently consume. As a result, spirally sliced hams are commonly sold at retail in half portions. The full ham is cut transversely through the femur in order to cut the product into a butt half and a shank half. The half that is adjacent to the butt end includes part of the femur and the aitch bone, while the half at the shank end of the product includes the rest of the femur and the shank bone. Both halves are commercially appealing in that only the ends of the femur are visible and the product is not unsightly. Also, the meat remains attached to the bone continuously completely around it, so slices are not prone to falling off or becoming deformed during subsequent handling and packaging.
Although spirally sliced ham products of this type have been well received, the amount of meat they contain is still excessive for smaller households and smaller occasions where the product is consumed by only a relatively small number of people. Consequently, there is the need to cut spirally sliced hams into smaller products in order to allow consumers the option of purchasing a smaller amount of meat. However, cutting the ham into smaller products presents significant problems which have not been overcome and such smaller spirally sliced ham products have not been commercially available.

Spirally sliced ham product and method and apparatus for production thereof


A switching power converter with a controlled startup mechanism includes a switching stage which provides a voltage Vout at an output node in response to a switching control signal, with the output node adapted for connection to a non-linear load. A feedback network compares a signal which varies with the current conducted by the load (Iload) with a reference signal, and provides the switching control signal so as to maintain Iload at a desired value. A capacitor connected to the output node provides a current Ic to the feedback network which varies with dVout/dt. The feedback network is arranged to limit dVout/dt in response to current Ic when Iload is substantially zero. In this way, large inrush currents or damage that might otherwise occur during startup are avoided.

1. Field of the Invention
This invention relates generally to switching power converters, and more particularly, to a switching converter with a startup mechanism which is useful when the load being driven is non-linear.
2. Description of the Related Art
Switching power converters can be configured using a number of different topologies, can be controlled with a wide variety of control methods, and can be used to drive different types of loads. One possible arrangement is shown in FIG. 1. In this example, a boost converter topology is employed, in which an inductor 10 is connected between an input voltage Vin and a switching element—here, a transistor 12; a diode 14 is connected to the junction of the inductor and transistor, and the converter's output voltage Vout appears on the cathode side of diode 14. Vout drives a load 16—here, a diode—connected between Vout and a node 20, and a resistor 22 having a resistance R1 is connected between node 20 and a common potential. A filter capacitor 24 is typically connected across load 16. When so arranged, the current Iload in load 16 can be controlled by controlling the voltage at node 20.
The voltage at node 20 is controlled by means of a feedback circuit. An error amplifier 26 receives a reference voltage Vref at one input and is connected to node 20 at a second input. The output of the amplifier drives a controller 28 which operates transistor 12. In operation, controller 28 operates transistor 12 as needed to make Vout the value needed to force the voltage at node 20 to Vref, thus causing Iload to be given by Vref/R1.
However, a converter arranged as shown can encounter a problem during the "startup" phase of its operation. When Vout is beginning to increase from zero, but is still below the forward voltage of diode 16, little to no current flows in resistor 22 and node 20 will be nearly equal to the common potential. With a voltage of zero at node 20, error amplifier 26 sees a large error and drives the controller to operate transistor 12 at a maximum duty ratio not controlled by feedback. This causes the average inductor current to rise rapidly, possibly to a level at which it saturates inductor 10 and possibly damages transistor 12.

Switching power converter with controlled startup mechanism


A method and apparatus for mutual authentication of a first and a second node is provided herein. During operation the first node sends a first authentication message to the second node comprising information needed for the second node to authenticate the first node and information needed by the second node for the second node to be authenticated by the first node. A second authentication message is received from the second node comprising information needed by the first node for the first node to be authenticated by the second node and information needed for the first node to authenticate the second node. Because the first and the second authentication messages comprise information needed for both the first and the second node to authenticate each other, messaging within the network is greatly reduced.

In many networks a node will need to authenticate another node prior to communication taking place between the nodes. Such authentication, may take place by performing a full authentication exchange with a network entity such as an Authentication, Authorization, Accounting server (AAA server) or an Extensible Authentication protocol server (EAP server) that is aware of the node's rights with respect to network access. Such authentication can be done through a variety of methods and generally involves much messaging between the two nodes and the server.
When networks require mutual authentication (MA), each node must authenticate the other node prior to communication taking place. In other words, when a first node wishes to communicate with a second node, the first node must authenticate the second node, while the second node must similarly authenticate the first node. Because authentication methods generally involve many messages between the two nodes and the server, and because two nodes will have to authenticate each other when performing mutual authentication, mutual authentication requires extensive messaging between both nodes and the server. Because of this, a need exists for a method and apparatus for performing mutual authentication within a network that reduces the messaging between the two nodes.

Method and apparatus for performing mutual authentication within a network


Apparatus to vary the effective spring rate of diaphragm actuators are described. An example spring assembly for use in a diaphragm actuator includes a first spring operatively coupled to a diaphragm plate and an actuator stem, and a second spring operatively coupled to the actuator stem and configured to be serially operative with the first spring. Additionally, a movable spring seat is disposed between the first and second springs.

Control valves are used in process control systems to manipulate a flowing fluid, to regulate a process variable to a desired set point, to convey or deliver a fluid from a source to a destination, etc. A control valve assembly typically includes a valve body, a shaft or stem, and an actuator to provide the motive power via the shaft or stem to operate the valve or position a plug or flow control member within the valve. A common type of actuator for use with a control valve assembly is a spring and diaphragm pneumatic actuator, which is commonly referred to as diaphragm actuator.
Typically, diaphragm actuators have a casing to house a diaphragm, a diaphragm plate, an actuator stem, and a spring assembly having one or more springs. The spring assembly applies a force against the diaphragm plate to return the actuator stem and a valve or other operator coupled to the stem to a known position in the absence of a control pressure applied to the diaphragm. In the case where a single spring is used to implement the spring assembly, the spring is typically centrally located on the diaphragm plate. Where multiple springs are used, the springs are typically distributed circumferentially about the center of the diaphragm plate and apply their respective forces directly to the diaphragm plate.
Regardless of whether one or multiple springs are used to implement the spring assembly, diaphragm actuators receive a control pressure to vary gas (e.g., air) pressure on one side of the diaphragm to move or stroke the actuator stem and thereby open and close or modulate a control valve that may be coupled to the actuator stem. The amount of control pressure required to move and maintain the actuator stem and, thus, the valve or other operator controlled by the actuator, at a given position along its range of stroke, typically equals the force exerted by the spring assembly plus the force exerted by the valve or other operator on the actuator stem. The force exerted by the spring assembly typically increases in a substantially proportional and linear manner as the actuator stem moves toward its fully stroked condition.
Additionally, the force exerted against the actuator stem by the valve or other operator may also vary (e.g., increase or decrease) as the position of the valve (e.g., the position of a valve plug, disc, etc.) changes from its zero stroke to its fully stroked condition. Such changes in the force exerted by the valve stem may be due to the flow characteristics of the valve, the pressure and flow rate of the fluid being controlled by the valve, etc. and, thus, may be substantially nonlinear. Thus, the amount of control pressure required to fully stroke the actuator and the valve or other operator coupled to the actuator may be relatively high and, as a result, may require the diaphragm, the casing, and other actuator components to withstand relatively high pressures. Further, because the force exerted by the actuator spring assembly typically varies linearly over the range of the actuator stroke and because the force exerted by the valve on the actuator stem may vary non-linearly or in some other manner, the position of the valve may vary in an undesirable manner relative to control pressure (e.g., substantially non-linearly).

Apparatus to vary effective spring rate for use in diaphragm actuators


A computer can be built into a pre-existing exercise machine so that the person using the machine can view the screen on the computer monitor while engaged in aerobic exercise activity. Keyboard mechanisms are secured to hand rails of the machine, whereby the person can operate the keyboards without releasing his grip on the hand rails. The computer monitor is adjustably supported to provide persons of different heights and head positionments a good view of the viewing screen on the monitor.

SUMMARY OF THE PRESENT INVENTION
This invention relates to an exercise machine and particularly to an exercise machine equipped with a built-in computer, whereby the person using the exercise machine can at the same time scan the computer monitor for information purposes or amusement purposes. By hooking the computer into the internet, the system will offer the person using the machine a vast amount of information.
Over the years, corporations have made substantial investments in fitness facilities for their employees. It has been shown that employee involvement in corporate fitness programs increases health, productivity and self-esteem, while at the same time reducing the health care costs of the corporation.
My invention is designed to increase the productivity of an employee while he is exercising his body on an exercise machine. The invention eliminates the problem of an employee failing to exercise because he has too much work to accomplish within the time available. With appropriate computer hook-ups, the employee will be able to access corporate voice, video and data networks within the corporate communication system, while at the same time enjoy the benefits of aerobic exercise on a stationary exercise machine. The invention can be applied to various types of exercise machines, e.g. treadmills, cycles and stepping machines, etc.
In addition to usage in the corporate environment, the invention can also be used in the commercial fitness (health) area and home fitness area. The individual using an exercise (fitness) machine will be able to surf the internet web, and enjoy an endless quantity of online news, information and entertainment features. The discomfort associated with exercise will be reduced, due to the fact that the individual can be engrossed in processing information of interest.
The nature of the invention will become more apparent from the attached drawings and description of an illustrative embodiment of the invention.
In summary, and in accordance with the above discussion, the foregoing objectives are achieved in the following embodiments.
1. In combination, an exercise machine having a front end and a rear end; a hollow post structure at the machine front end, an exercise function control panel atop said hollow post structure, and two hand rails at least partially supported by said post structure;
a television monitor supported by said post structure above said hand rails; and PA1 keyboard means for controlling the pictoral display on said monitor; said keyboard means comprising an array of thumb actuated switches supported by one of said hand rails, whereby the person using the exercise machine can operate said keyboard means while grasping said hand rails.
2. The combination, as described in paragraph 1, and further comprising an adjustable support means for said television monitor; said adjustable support means comprising a support arm extending from said post structure, and a television monitor platform having a hinged connection to said support arm.
3. The combination, as described in paragraph 1, wherein said television monitor has a viewing screen facing the occupant space above the exercise machine; and an adjustable support means for said television monitor; said adjustable support means comprising a support arm extending from said post structure, and a television platform having a hinged connection to said support arm; said hinged connection having a horizontal hinge axis that enables the television monitor to tilt around the hinged connection, whereby the viewing screen is readily viewable by different persons using the exercise machine.
4. The combination, as described in paragraph 1, wherein said keyboard means comprises an elongated switch panel angling downwardly from said one hand rail.
5. The combination, as described in paragraph 1, wherein said keyboard means comprises an elongated switch panel angling downwardly from said one hand rail, and means for varying the angulation of said switch panel relative to said one hand rail.
6. The combination, as described in paragraph 1, wherein said keyboard means comprises an elongated switch panel angling downwardly from said one hand rail; said switch panel comprising a printed circuit board, a flexible non-conductive cover sheet overlying said circuit board, and a series of electrical contacts spaced along said cover sheet for selective engagement with said printed circuit board in response to thumb pressure on said cover sheet.
7. The combination, as described in paragraph 1, wherein said exercise machine is a treadmill.
8. The combination, as described in paragraph 1, and further comprising a hollow support arm extending from said post structure in supporting relationship to said television monitor, and a microprocessor located in said hollow support arm.
9. The combination, as described in paragraph 8, and further comprising first lead wiring extending within said one hand rail for operatively connecting said keyboard means to said microprocessor, and second lead wiring linking said microprocessor to said television monitor.

Exercise machine having a built-in computer


The present invention provides a technique for converting a multi-threaded application configured to execute on a uniprocessor (UP) system to one that executes on a multiprocessor (MP) system. Unlike previous approaches, a novel scheduling technique is employed so that different UP-coded user-level threads (“sthreads”) can execute concurrently in the MP system without having to rewrite their original code. To that end, the UP-coded sthreads are organized into different concurrency groups, each of which defines a set of one or more sthreads not permitted to execute concurrently. By grouping the UP-coded sthreads in this manner, different concurrency groups can be scheduled to execute their UP-coded sthreads at substantially the same time without incorporating traditional synchronization mechanisms into the sthreads' original UP code.

A computer system comprising a single processor is referred to as a uniprocessor (UP) system. Typically, the UP system initially loads an operating system to manage all other program code executed by the processor. A set of core services provided by the operating system is often coded in the operating system's kernel. Such services may include, among other things, providing file-system semantics, input/output (I/O) operations, memory management, network communications, and the like. A kernel may be embodied as a microkernel that organizes the kernel's services in relatively small, configurable is modules. In this sense, the Data ONTAP™ operating system, available from Network Appliance, Inc. of Sunnyvale, Calif., is an example of an operating system implemented as a microkernel.
In general, kernel code in a uniprocessor system may be divided into multiple segments of program code. Each program defines a specific set of ordered operations that are performed by the processor. A process refers to an instance of a program being executed by the processor. Accordingly, UP system resources, such as memory address space and processor bandwidth, may be divided among multiple processes. Moreover, the resources may be allocated to both user and kernel processes. Here, a kernel process implements one or more of the kernel's services. In contrast, a user process provides a service, such as word processing, web browsing, image editing, software developing, etc., for a user.
Typically, a process scheduler controls when different processes are executed by the processor. In some implementations, the process scheduler is integrated into the kernel ("kernel-level"). Other implementations code the scheduler outside the kernel ("user-level"). The process scheduler maintains a run queue that governs the order in which processes are executed by the processor. Thus, the scheduler dequeues a process for the processor to execute, then re-enqueues the process after an appropriate time interval elapses or the processor finishes the process's execution.
In practice, a process in a UP system often may be interrupted (i.e., "stalled") by the processor so other operations can be performed in the system. For instance, the kernel may recognize an I/O interrupt signal having higher priority than the executing process. In this case, the process may be stalled while the processor performs appropriate I/O operations. Similarly, other process interruptions may include memory page faults, segmentation faults, error handling routines, and so forth.
A process often employs one or more threads to minimize processor stalling and make more efficient use of the processor's bandwidth while the process is being executed. Each thread is an instance of a sequence of the process's program code. For instance, a user-level thread is an instance of a sequence of the process's program code which is not included in the kernel's code. The threads allow the process to keep the processor busy, even in the presence of conventional process interruptions, such as I/O interrupts, page faults, etc. To that end, the process can perform "blocking" operations when an executing thread is interrupted, thereby suspending that thread's execution and beginning execution of a different thread. In this manner, the process can continuously execute thread code for as long as the process is scheduled to execute on the processor.
As used herein, one or more workloads may be multiplexed on an execution vehicle that executes the workloads in an orderly manner. For example, when the workloads are user-level threads, one or more of the user-level threads may be multiplexed on an execution vehicle, such as a kernel process, through which they are executed. Specifically, the process scheduler may allot a time interval for a kernel process to execute, and the kernel process may divide its allotted time among its user-level threads. To that end, a thread scheduler may be used to control when the process's different user-level threads are executed by the processor. The thread scheduler allows the process to continue executing even when one or more threads block, thereby maximizing utilization of the process's scheduled processor during the process's time slice. Thus, in the absence of other processors, the thread scheduler maximizes the system's processor utilization.
Each thread is associated with a thread context that supplies the processor with meta-data information needed to execute the thread. For instance, a thread's associated program-counter value, stack-pointer value, hardware register values, etc. are stored in its thread context. Each thread context is typically stored in a corresponding data structure called a thread control block (TCB). Although threads are associated with their own thread-specific meta-data information, multiple threads may share access to a common process context that contains process-specific meta-data information shared among the threads. The process context may include global variables, memory address spaces, semaphores, etc., defining the process's computing environment. Typically, the process context is stored in a corresponding data structure known as a process control block (PCB).
Thread contexts are typically smaller (i.e., "lightweight") than their associated process contexts, so switching thread contexts generally consumes less resources, such as processor bandwidth and memory, than switching process contexts. Furthermore, a new thread is typically not allocated many of the resources, such as "heap" memory, run-time libraries, etc., that are allocated to a new process. In fact, only a TCB and a call stack are typically allocated when a new thread is created; other resources are previously allocated at the time the thread's process is created. Thus, threads may be dynamically created and destroyed using less overhead (i.e., resource consumption) than processes since they not only require less context information, but also require less resource allocation. Notably, threads need not be lighter than processes in every dimension. For example, it is possible to have threads which consume less memory than a process, even if they execute less efficiently and/or require more time to context switch.
While multi-threading techniques may be used to ensure efficient use of the processor in a UP system, threads also facilitate efficient processing in a multi-processor (MP) system. As used herein, two or more processors in an MP system execute concurrently if they are scheduled to execute program code at substantially the same time. A process scheduler in an MP system traditionally determines which processors execute processes in the scheduler's run queue. For instance, the scheduler may assign a first processor to execute a first process, and at the same time assign a second processor to execute a second process. In this case, the first and second processors may concurrently execute user-level threads multiplexed over their respective first and second processes. Accordingly, thread schedulers associated with the first and second processes may be used to schedule user-level thread execution on the first and second processors, thereby maximizing utilization of both processors' bandwidths.
However, problems often arise in an MP system when concurrently executing processors attempt to access (i.e., read or write) the same data at substantially the same time. For example, a processor may attempt to modify data, while the same data is being accessed or processed by another processor. As a result, modifications made to the data may be known to one processor and not the other, thus causing inaccuracies or errors in one or both of the processors' execution. To prevent such conflicts, synchronization mechanisms are often employed to give processors exclusive access to shared data. Conventional synchronization techniques include, inter alia, global locks, monitors, domains and master/slave approaches.
A unique global lock may be associated with each data structure allocated in an MP system. According to this synchronization technique, a sequence of program code, such as a thread or process, must "acquire" a global lock before it can access data in the lock's associated data structure. When the program code is finished accessing the data, it "releases" the global lock. For example, a thread may acquire a global lock by setting the lock to a first value (e.g., "1"), then later release the lock by setting it to a second value (e.g., "0"). Illustratively, the lock's value may be stored in a special "lock" data structure, although those skilled in the art will appreciate other implementations are possible as well. Broadly stated, a global lock is a type of mutual exclusion ("mutex") lock since it can only be possessed by one thread or process at a time. Threads or processes may either "block" while waiting for mutex locks to be released or may loop to repeatedly check the status of a lock held by another thread or process. The latter process is called "spinning" and mutex locks that require spinning are known as spin locks.
While a global lock synchronizes access to an individual data structure, a monitor synchronizes execution of different segments of program code, such as threads. In general, a monitor is a program that manages the interactions of other segments of program code. Typically, a monitor is associated with a monitor lock that can be acquired by only one thread or process at a time. For example, suppose one or more data structures is shared among a set of threads, and each of the threads executes through a common monitor. Since the monitor's lock only allows one of the threads to acquire the monitor at any given time, the lock prevents multiple threads from accessing the shared data structures at the same time. When the threads are scheduled to execute on different processors in an MP system, the monitor can therefore control their concurrent access to the shared data structures.
A domain typically protects a larger amount of program code than a monitor. More specifically, a programmer may organize a program's code into a plurality of distinct code segments, or domains, that can execute concurrently in an MP system. However, processes and threads within the same domain usually may not execute concurrently. Further, each instance of a function or data type allocated in a domain is accessible to only one specific thread or process executing in that domain. Each domain is typically associated with its own domain-specific context information and resources. Before a domain can load its context information and execute its processes and threads, the domain first must acquire its associated domain lock. The domain lock prevents multiple instances of the domain's threads and processes from executing concurrently on different processors in the MP system. Synchronization using domains is usually determined at compile-time and is further described in more detail in commonly assigned application Ser. No. 09/828/271, entitled "Symmetric Multiprocessor Synchronization Using Migrating Scheduling Domains," to Rajan et al., which is hereby incorporated by reference as though fully set forth herein.
A master/slave implementation assigns each process that can execute in an MP system to a "master" or "slave" processor. Specifically, program code that cannot execute concurrently in an MP system is assigned to a master processor. Other program code may be executed concurrently in the MP system by one or more slave processors. In this manner, a software programmer controls which threads and processes execute concurrently in the MP system by choosing whether each process or thread executes on a master or slave processor. Thus, unlike the other synchronization techniques previously described, the master/slave approach requires the programmer to predetermine which processors will execute different processes in a software application.
While conventional synchronization techniques enable multiple processes and threads to execute concurrently in an MP system, they are not easily implemented in program code originally designed to execute in a UP system. For instance, converting a UP-coded thread to employ global locks in an MP system requires adding extra code to the thread for acquiring and releasing the locks. Firstly, this additional code may be time consuming to incorporate. Secondly, the extra code adds overhead to the thread even if the thread never executes concurrently with other threads, and therefore may unnecessarily slow the thread's execution. Thirdly, incorporating this additional code into the UP-coded thread may result in the thread consuming an unacceptable amount of memory resources in the MP system. In addition, MP-system synchronization using global locks may lead to "deadlocks." For example, suppose a first thread in the MP system seeks to acquire global lock A then B, and a second thread seeks to acquire global lock B then A. In this case, the first and second threads block one another. If one of the locks A or B is critical to the system, the entire system can freeze. Programming around such issues can considerably complicate MP development and reduce the efficiency of the MP code.
Similarly, converting a UP-coded thread to employ monitors in an MP system also requires additional code to be added to the thread. For example, the thread must be configured to acquire and release monitor locks. Further, the thread must be modified to incorporate signaling and blocking mechanisms typically associated with monitors. In sum, the UP-coded thread would have to be re-written within a monitor paradigm. This process of re-writing and redesigning the thread's code is costly in terms of both time and development resources.
In general, some system architectures align well with domain synchronization techniques, while others do not. Thus, converting a UP-coded thread to implement domains in an MP system may require re-architecturing of the thread's code to ensure its compatibility with the MP system's architecture. For instance, a programmer typically must divide the original UP-coded program code into individual domains. Although statically allocating blocks of domain code may simplify debugging operations, the process of splitting large amounts of program code into domains is typically an arduous process. For instance, the UP-coded processes and threads within a domain must be modified in accordance with the domain's synchronization mechanisms, e.g., to properly utilize the domain's associated domain lock, domain context, domain identifiers, etc. The process of re-writing and redesigning the UP threads' code to employ domain synchronization mechanisms typically consumes excessive software-development time and resources. Furthermore, domains may be of limited use in the MP system since each instance of a function or data type allocated in a domain is typically only manipulated by a single thread or process.
Lastly, converting a UP-coded thread to execute in a master/slave configuration requires a programmer to assign the thread to execute on a master or slave processor in the MP system. According to the master/slave approach to concurrency, a UP-coded thread assigned to a slave processor must be rewritten to make it MP safe with respect to all of the other threads in the MP system. Thus, the master/slave approach suffers the disadvantage of having to incorporate conventional synchronization mechanisms into UP thread code assigned to the slave processors.
It is therefore desirable to convert a UP-coded thread for use in an MP system without having to modify the original thread code. The thread should not have to be recoded or augmented with additional code in order to execute in the MP system. Also, the thread should be able to cooperate with other threads and/or processes to access instances of functions and data types allocated in the MP system. Further, the UP thread code should not be "bound" to any particular hardware implementation, and should easily scale when processors are added or removed in the MP system. It is also desirable that the same thread code can function in both the UP and MP systems.

Technique for dynamically restricting thread concurrency without rewriting thread code


Inorganic semiconductive films are made by depositing a suitable precursor substance upon a substrate, irradiating the precursor substance with electromagnetic radiation to form a nascent film, and heating the nascent film at a predetermined temperature to form an inorganic semiconductive film.

Electronic devices such as field-effect transistors (FET's) are conventionally formed on a substrate, such as a silicon semiconductor substrate or an insulating substrate, and are combined with interconnections to form integrated circuits, which may be used, for example, as an array of transistors to address pixel elements in displays. The integrated assemblies have typically been built up by forming a number of layers, each layer being formed by depositing a film of semiconductor, insulator or electrical conductor material and by subsequently patterning the film of deposited material, typically by photolithography in a subtractive process (i.e., a portion of a blanket-deposited film is selectively removed). Conventionally, fabrication of semiconductor devices uses a number of different processes performed sequentially, and such sequential methods can increase the time, and therefore the cost, of device fabrication.
Inorganic semiconductive films have been formed by many known methods, including deposition by evaporation, chemical vapor deposition, physical vapor deposition (i.e., sputtering), etc., and subsequent patterning by photolithography.
Inkjet printing methods have been used for many purposes, including ordinary printing of ink on paper, formation of 3-dimensional objects, deposition of photopolymerizable solutions for photolithographic patterning, printing of color filters for liquid-crystal displays, printing of light-emitting layers for electroluminescent devices, printing of conductive traces on circuit boards and integrated circuits, and many other purposes.
Organic semiconductor films have been inkjet printed, and these materials are suitable for some applications, but they have not provided carrier mobilities (i.e., for current drive capability) achievable with inorganic semiconductor films. Another well-known disadvantage of organic semiconductors is that the performance (e.g., on-to-off current ratio) strongly declines over a relatively short period of time, which precludes their use for long-term-use products (e.g., two years or more of use).
Inorganic thin-film semiconductors of polycrystalline silicon and amorphous silicon have many attractive qualities, including lifetime and performance qualities, but patterned films of polycrystalline silicon or amorphous silicon have not been prepared successfully by any additive process.
Thus, there is a need for efficient and economical methods for forming patterned inorganic semiconductive films.

Methods for making and using high-mobility inorganic semiconductive films

Hewlett-Packard Development Company, L.P.


An absorbent article, such as diapers and the like, comprises a liquid pervious topsheet, a liquid impervious backsheet joined with said topsheet and an absorbent core positioned between said topsheet and said backsheet, said backsheet comprising a flexible starch based film comprising starch, a synthetic thermoplastic polymer of at least one ethylenically insaturated monomer, said polymer having repeating units provided with at least a polar group, wherein the starch and the polymer form an at least partially interpenetrated network, and moisture.

A wide variety of absorbent articles designed to be efficient for the absorption of body fluids such as blood, urine, menses, and the like, are known. Disposable products of this type generally comprise some sort of fluid-permeable topsheet material, an absorbent core, and a fluid-impermeable backsheet material.
Heretofore, such absorbent structures have been prepared using, for example, topsheet materials prepared from woven, nonwoven, or porous formed-film polyethylene or polypropylene materials. Backsheet materials typically comprise flexible polyethylene sheets. Absorbent core materials typically comprise wood pulp fibers or wood pulp fibers in combination with absorbent gelling materials.
One aspect of such absorbent articles that has recently been considered is their disposability. Although such products largely comprise materials which would be expected ultimately to degrade, and although products of this type contribute only a very small percentage of the total solid waste materials generated by consumers each year, nevertheless, there is currently a perceived need to devise such disposable products from materials which are compostable.
A conventional disposable absorbent product is already to a large extent compostable. A typical disposable diaper, for example, consists of about 80% of compostable materials, e.g., wood pulp fibers, and the like. In the composting process soiled disposable absorbent articles are shredded and commingled with organic waste prior to the composting per se. After composting is complete the non-compostable particles are screened out. In this manner even today's absorbent articles can successfully be processed in commercial composting plants. Nevertheless, there is a need for reducing the amount of non-compostable materials in disposable absorbent articles. There is a particular need to replace polyethylene backsheets in absorbent articles with liquid impervious films of compostable material, because the backsheet is typically the largest non-compostable component of a conventional disposable absorbent article.
It is, therefore, an object of the present invention to provide absorbent articles having a liquid impervious backsheet comprising a compostable polymer.


A moving plate (30) has a substantially rectangular bottom portion (38) and a surrounding wall (32) projecting from the outer periphery of the bottom portion (38). Shake preventing projections (37) are formed on the opposite outer side surfaces of the surrounding wall (32) including the shorter sides of the bottom portion (38) and extend from an upper edge of the surrounding wall (32) over the entire vertical length. The height of the shake preventing projections (37) exceeds a corresponding clearance dimension between the surrounding wall (32) and a receptacle (11). A lever (20) is used to realize a smooth connecting operation by preventing the moving plate (30) from being inclined while two housings (10, 40) are being connected.

1. Field of the Invention
The invention relates to a connector and to a connector assembly.
2. Description of the Related Art
U.S. Pat. No. 5,269,696 discloses a multi-contact connector with a male housing and a moving plate movably disposed in the male housing. Male terminal fittings are mounted in male housing, and tabs of the male terminal fittings penetrate through the moving plate. Thus, the moving plate protects and positions the tabs.
A conventional moving plate is smaller than the receptacle by a specified amount to ensure that the moving plate can be inserted easily into the receptacle. However, this clearance can cause the moving plate to incline and move less smoothly. A smaller clearance will make the moving plate less likely to incline, but can complicate the insertion of the moving plate into the receptacle.
The invention was developed in view of this problem, and an object thereof is to ensure easy insertion of the moving plate without inclination.


This invention discloses a semiconductor device comprising a semiconductor substrate, a first conducting layer formed on the surface of the semiconductor substrate, an insulating layer formed above the semiconductor substrate, the insulating layer having a contact hole reaching the first conducting layer to expose it, a second conducting layer formed on the insulating layer, the sidewall of the contact hole, and the first conducting layer, and an anti-oxidation layer formed on at least part of the surface of the second conducting layer.

1. Field of the Invention
This invention relates to a semiconductor device and a method of manufacturing the same, and more particularly to a high-density LSI with microscopic interconnections.
2. Description of the Related Art
In conventional semiconductor devices, for example, DRAM memory cells with bit lines have the structure as shown in FIGS. 1A through 1C. FIG. 1B is a sectional view taken along line 1B--1B of FIG. 1A, and FIG. 1C is a sectional view taken along fine 1C--1C of FIG. 1A. Numeral 701 indicates a p-type semiconductor substrate, 702 an element isolating oxide film, 703 a silicon oxide film, 704 the gate electrode of a MOS FET, 705 an n-type diffused layer, 706A and 706B interlayer insulating films, 707A and 707B contact holes, 708 the bottom electrode of a capacitor, 709 a capacitor insulating film, 710 the top electrode of the capacitor, and 711 a bit line.
On the memory cell, a flattened interlayer insulating film, such as a BPSG film, is generally formed. Formed on the insulating film are metal interconnections (of, for example, aluminum), on which a passivating film is formed to complete the DRAM. Referring to FIGS. 1A through 3B, a manufacturing method of DRAMs will be explained hereinafter, covering all the way to the final products.
As shown in FIG. 1B, the element isolating oxide film 702 is first formed on the p-type semiconductor substrate 701. On an element region of the p-type semi-conductor substrate 701, the silicon oxide film 703, gate electrode 704, and n-type diffused layer 705 are each formed in a known manner to complete the MOS FET. After the interlayer insulating film 706A is formed all over the surface, the contact hole 707A leading to the n-type diffused layer (source) 705 is made. At the contact hole 707A, the bottom electrode 708 of the capacitor, capacitor insulating film 709, and the top electrode 710 of the capacitor are formed to produce a cell capacitor of the DRAM. Over the whole surface, the interlayer insulating film 706B is formed, and then the contact hole 707B reaching the n-type diffused layer (drain) 705 is made. After this, on the interlayer insulating film 706B and contact hole 707B, the bit line 711 is formed by sputtering with silicide such as MoSi.sub.2 or WSi.sub.2.
Then, as shown in FIG. 2A, over the entire surface, the interlayer insulating film 712 such as silicate glass (BPSG film) containing, for example, boron (B) or phosphorus (P) is formed. The interlayer insulating film 712 has larger steps in the vicinity of the bit line 711 as shown in FIG. 2B.
Next, as shown in FIGS. 3A and 3B, high-temperature heat treatment (annealing) is performed in an oxidizing atmosphere to flatten the surface of the interlayer insulating film 712. At this time, an oxidizer (e.g., oxygen) penetrates the interlayer insulating film 712 and oxidizes the bit line 711 made of a silicide film. As a result, on the surface of the bit line 711, the oxide film 713 is formed. The metal interconnections (made of, for example, Al) 714 are formed on the flattened interlayer insulating film 712. Over the entire surface, the passivating film 715 is formed to complete the DRAM.
With the aforementioned manufacturing method, however, it is known that the bit line 711 formed by sputtering has a poor step coverage, with the thickness of the bit line 711 being thinner on the sidewall of the contact hole 707 than on the top surface. In this state, when heat treatment is carried out in an oxidizing atmosphere, the thinner portions of the bit line 711 on the sidewall of the contact hole 707B are all oxidized, leading to breaks in interconnections and an increase in the resistance. Because of this, neither sufficient yield nor high reliability can be achieved.
As noted above, conventional semiconductor devices have a poor step coverage of interconnections at the contact holes. Because of this, a subsequent heat treatment causes the oxidizer to react upon the interconnections to form an oxide film, which results in breaks and increase in the resistance at the thinner portions of the interconnections in the contact holes.

Semiconductor device and a method of manufacturing same


A non-volatile memory device includes an array of non-volatile memory cells. A memory cell in the array of memory cells includes a first resistive element including a first terminal and a second terminal, a second resistive element including a first terminal and a second terminal, and a select transistor including a gate electrode coupled to a word line, a first current electrode coupled to the first terminal of the first resistive element and the first terminal of the second resistive element, and a second current electrode coupled to a bit line. The second terminal of the first resistive element is coupled to a first source line, and the second terminal of the second resistive element is coupled to a second source line.

Field
This disclosure relates generally to integrated circuits, and more particularly, to non-volatile random access memories (NVRAMs).
Related Art
Non-volatile memories (NVMs) have become very important in a variety of applications but typically have characteristics that make them difficult to use as a random access memory. Some of the difficulties are very slow program and erase times, inability to erase one bit at a time, and high voltage requirements for program and erase. Some of the resistive NVMs, such as magnetoresistive tunnel junctions (MTJs) overcome these difficulties making them a candidate for use as a random access memory (RAM). Other difficulties such as large space requirements make this less attractive.
Accordingly there is a need to provide further improvement in obtaining NVRAMs.

Non-volatile random access memory (NVRAM)


The invention provides a composite catalyst containing a first component and a second component. The first component contains nanosized gold particles. The second component contains nanosized platinum group metals. The composite catalyst is useful for catalyzing the oxidation of carbon monoxide, hydrocarbons, oxides of nitrogen, and other pollutants at low temperatures.

Removing the toxic pollutants in automotive exhaust has been an intense focus of the automotive industry over the last several decades. In particular, the emissions regulations for fuel-efficient diesel engines that were implemented in 2007 and 2010 have resulted in a new generation of emissions control technologies. One of these technologies that has a multifunctional purpose is the diesel oxidation catalyst (DOC). Its function is primarily to oxidize carbon monoxide (CO) and unburned hydrocarbons (HCs) to carbon dioxide (CO2) and water (H2O), but also is used to oxidize a portion of the nitric oxide (NO) to nitrite (NO2) for use in low temperature soot oxidation in the diesel particulate filter (DPF) and to enable the optimal NO/NO2 ratio for the selective catalytic reduction (SCR) of NOx with NH3 over metal-zeolites.
The active component of the DOC is typically Pt and/or Pd. These catalysts usually reach 90% conversion between 200° C. and 350° C., and consequently, the catalysts are not active under "cold-start" or low load/idling conditions. The temperature limitation is problematic since more than 50% of the emissions from an engine occur in the first 2-3 minutes after cold start. Thus, as emissions regulations become more stringent, meeting the emission regulations will require increased activity during this warm-up period. To further complicate matters, the increased Corporate Average Fuel Economy (CAFE) standards that will be implemented over the next decade will result in the introduction of more fuel-efficient engines. These engines will have even lower exhaust temperatures, which further necessitates the need for increased emissions control activity at low temperatures. Higher Pt/Pd loadings may help to increase the catalytic efficiency, but such methods would be too expensive and are more subject to particle sintering which degrades performance over time.
Other options to meet the emissions standards include hydrocarbon/NOx absorbers which operate by trapping pollutants below temperatures at which catalysts would treat them and, subsequently, releasing the pollutants once temperatures rise to a point at which a downstream catalyst can effectively treat them. However, while this pathway would help mitigate emissions at low exhaust temperatures, additional complexity and cost are drawbacks to implementation.
Thus, oxidation catalysts in emissions control systems, such as automotive exhaust that generally contain Pd/Pt require exhaust temperatures above 200° C. to operate; and under low-temperature conditions, oxidation of CO and hydrocarbons are challenging. As engine efficiency improves and exhaust temperature decreases, there is an increasing demand for high emissions control performance at low temperatures. Therefore, it becomes imperative to design new catalysts that are active at low operating temperatures.


Suggested are biocide compositions, comprising (a) 2-Furoic and/or Tetrahydro-2-furoic acid amides; (b) Biocides and optionally (c) Oil components or co-solvents and/or (d) Emulsifiers. The compositions show high solubility for biocides and improved emulsification properties.

Biocides, and in particular pesticides such as fungicides, insecticides and herbicides, are important auxiliary agents for agriculture in order to protect and to increase crops yields. Depending on the various and often very specific needs a magnitude of actives exist which show very different chemical structures and behaviors. Nevertheless, it is well known from the state of the art that it remains difficult to prepare concentrates of these actives which are exhibiting a satisfying stability, especially if stored at very low or elevated temperatures over a longer period. As a matter of fact, the solutions show a strong tendency to either separate or to form crystals, which makes it necessary to re-disperse the actives in the compositions prior to every application in order to obtain a homogenous product. Due to the fact that in spray equipments, which are customarily used for the application of aqueous formulations of plant treatment agents, several filters and nozzles are present, an additional problem appears which is related to the blocking of these filters and nozzles as a result of crystallizing active compound during the application of aqueous spray liquors based on solid active compounds.
The use of certain derivatives of lactic acid, in particular lactamides is well known from the state of the art. For example, German patent application DE 4112873 A1 discloses a composition comprising a fungicide (triflorine), lactic acid dimethyl amide and an emulsifier. The use of various lactamides as solvents in agriculture is described in WO 2006/040022 A2 (Bayer), WO 2007/107745 A2 (Syngenta), or EP 1961301 A2 (Cognis).
Products generally found in the market present either a limited content of biocides and/or hazardous labeling and a non favorable eco-toxicological profile.
The problem underlying the present invention has been to identify suitable new solvents for developing new biocide compositions with equal or higher contents of actives than obtainable in the market. The new solvents need to be safe and environmental friendly and should allow obtaining concentrated biocide compositions (on average more than 15% active matter) regardless of the chemical structure of the biocide, in other words they need to be excellent solvents for a wide range of herbicides, insecticides and fungicides. Finally, another object of the invention has been to design emulsifiable concentrates formulations with specific co-solvents and emulsifier system providing superior emulsion stability, in particular with respect to opacity and layering.

Biocide compositions comprising amides based on 2-furoic acid and its derivatives


An electronic accounting-document system, an electronic accounting-document processing method, an accounting-document creation unit and an accounting-document processing unit are provided in which the time that is taken to transmit and process accounting-document data can be shortened. An accounting-document creation section creates accounting-document data on an accounting document; an accounting-document data division section divides the accounting-document data into accounting-document present data and accounting-document absent data; a data transmission section transmits the accounting-document present data and the accounting-document absent data which is received by a data receiving section; an accounting-document data monitoring section identifies the accounting-document present data and the accounting-document absent data; a distribution accounting-document data creation section processes and stores, as accounting-document data for distribution, the accounting-document present data so that it can be output to a CD-R; and a distribution accounting-document data updating section updates the distribution accounting-document data, based on the accounting-document absent data.

1. Field of the Invention
The present invention relates to an electronic accounting-document system which allows an accounting-document processing unit to process accounting-document data that is created by an accounting-document creation unit. It also relates to an electronic accounting-document processing method, an accounting-document creation unit and an accounting-document processing unit for the same purpose.
2. Description of the Related Art
Conventionally, an electronic accounting-document system is known in which an accounting-document creation unit and an accounting-document processing unit are connected (e.g., refer to Japanese Unexamined Patent Publication No. 6-324820 specification). In that system, first, the accounting-document creation unit creates accounting-document data and transmits the accounting-document data to the accounting-document processing unit. Then, the accounting-document processing unit receives and manages the accounting-document data. In such an electronic accounting-document system, the recent computerization prompts the recording of accounting-document data in a CD-R.
In such an electronic accounting-document system as described above, an accounting-document creation unit transmits data, in each accounting-document unit, to an accounting-document processing unit, whether or not it includes accounting-document data. For example, in the case where a certain company handles accounting-document data, the total number of pieces of accounting-document data which actually includes accounting-document data is 18300. On the other hand, the total number of pieces of so-called 0-piece accounting-document data which does not include any accounting-document data is 102406. In that case, the accounting-document creation unit transmits blank data, some 100-thousand times, to the accounting-document processing unit, even though they have none of accounting-document data. This requires the accounting-document processing unit to process about 100-thousand pieces of blank accounting-document data.
In the above described conventional electronic accounting-document system, even if there is no accounting-document data, the accounting-document creation unit transmits data to the accounting-document processing unit. Hence, the time is wasted on such a transmission. Besides, in the above described conventional electronic accounting-document system, even though no accounting-document data exists, the accounting-document processing unit needs to process the data which has been transmitted from the accounting-document creation unit. This processing also takes a long time.

Electronic accounting-document system, electronic accounting-document processing method, accounting-document creation unit, and accounting-document processing unit