Durable asymmetric composite membranes consisting of a film of cross-linked poly(ether ether ketone) adhered to a sheet of hydrophilic microporous poly(ethylene) are disclosed. The membranes are suitable for use in the recovery or removal of water from feed streams where repeated clean-in-place protocols are required such as in the processing of dairy products.

It is well-known to use photografting to modify the surface of films, sheets and molded objects formed from polyolefins. For example, the publication of Tazuke and Kimura (1978) discloses photografting onto poly(propylene), poly(ethylene) and several other polymer films using benzophenone as a sensitizer. In this publication the choice of solvent and sensitizer was noted to be very important. The publication of Ang et al (1980) discloses an irradiation procedure where the sensitizer is dissolved in the hydrophilicitizing agent solution and can be used for the photosensitized copolymerization in high yields of styrene, 4-vinyl pyridine and methyl methacrylate to poly(propylene). Again, this publication notes that the reaction was found to be very specific to certain types of sensitizers.
The publication of Ogiwara et al (1981) discloses the photografting on poly(propylene) and low-density poly(ethylene) (LDPE) films on which sensitizers were coated beforehand. The sensitizers coated on films enabled vinyl hydrophilicitizing agents, such as methyl methacrylate, acrylic acid and methacrylic acid to graft easily with high yields. The hydrophilic hydrophilicitizing agents acrylic acid and methacrylic acid were conveniently grafted using them in aqueous solution in a liquid phase system. The publication of Allmer et al (1988) discloses the modification of surfaces of LDPE, high-density poly(ethylene) (HDPE) and polystyrene by grafting with acrylic acid. The grafting is performed in the vapor-phase and increased the wettability of the polymer. It was observed that acetone was able to initiate grafting and was found to promote and direct grafting to the surface. The publication of Allmer et al (1989) discloses the grafting of the surface of LDPE with glycidyl acrylate and glycidyl methacrylate by photoinitiation. Acetone and ethanol were used as solvents, with acetone yielding slightly more grafting at the surface.
The publications of Yao and Ranby (1990a, 1990b and 1990c) disclose inter alia a process for the continuous photoinitiated graft copolymerization of acrylamide and acrylic acid onto the surface of HDPE tape film. The process is performed under a nitrogen atmosphere using benzophenone as the photoinitiator. It was noted that pre-soaking was very important for efficient photographing within short irradiation times. The application of this pre-soaking photografting method to poly(ethylene terephthalate) (PET) was also disclosed. In this context acetone was found to be a somewhat better solvent than methylethyl ketone and methylpropyl ketone. When applied to a continuous process for the photochemically induced graft polymerization of acrylamide and acrylic acid of poly(propylene) (PP) fiber surface under a nitrogen atmosphere, optimal concentrations of hydrophilicitizing agent and initiator in the pre-soaking solution were determined.
The publications of Kubota and Hata (1990a and 1990b) disclose an investigation of the location of methacrylic acid chains introduced into poly(ethylene) film by liquid and vapor-phase photograftings and a comparative examination of the photografting behaviours of benzil, benzophenone and benzoin ethyl ether as sensitizers. In these latter studies poly(methacrylic acid) was grafted onto initiator-coated LDPE film.
The publication of Edge et al (1993) discloses the photochemical grafting of 2-hydroxyethyl methacrylate (HEMA) onto LDPE film. A solution phase method is used to produce a material with increased wettability. The publication of Singleton et al (1993) discloses a method of making a polymeric sheet wettable by aqueous solvents and useful as an electrode separator in an electrochemical device. The polymeric sheet is formed from fibers which comprise poly(propylene) alone and is distinguished from a membrane formed from a microporous polymer sheet. The publication of Zhang and Ranby (1993) discloses the photochemically induced graft copolymerisation of acrylamide onto the surface of PP film. Acetone was shown to be the best solvent among the three aliphatic ketones tested.
The publications of Yang and Ranby (1996a and 1996b) disclose factors affecting the photografting process, including the role of far UV radiation (200 to 300 nm). In these studies benzophenone was used as the photoinitiator and LDPE film as the substrate. Added water was shown to favour the photografting polymerisation of acrylic acid on the surface of polyolefins, but acetone was shown to have a negative effect due to the different solvation of poly(acrylic acid) (PAA).
The publication of Hirooka and Kawazu (1997) discloses alkaline separators prepared from unsaturated carboxylic acid grafted poly(ethylene)-poly(propylene) fiber sheets. Again, the sheets used as a substrate in these studies are distinguished from a membrane formed from a microporous polymer sheet.
The publication of Xu and Yang (2000) discloses a study on the mechanism of vapor-phase photografting of acrylic acid in LDPE. The publication of Shentu et al (2002) discloses a study of the factors, including the concentration of hydrophilicitizing agent, affecting photo-grafting on low-density LDPE. The publication of El Kholdi et al (2004) discloses a continuous process for the graft polymerisation of acrylic acid from hydrophilicitizing agent solutions in water onto LDPE. The publication of Bai et al (2011) discloses the preparation of a hot melt adhesive of grafted low-density poly(ethylene) (LDPE). The adhesive is prepared by surface UV photografting of acrylic acid onto the LDPE with benzophenone as the photoinitiator.
The publication of Choi et al (2001) states that graft polymerisation is considered as a general method for modifying the chemical and physical properties of polymer materials.
The publication of Choi (2002) discloses a method for producing an acrylic graft polymer on the surface of a polyolefin article comprising the steps of immersing the article in a solution of an initiator in a volatile solvent, allowing the solvent to evaporate, and then immersing the article in a solution of an acrylic hydrophilicitizing agent before subjecting the article to ultraviolet irradiation in air or an inert atmosphere. Acrylic acid is used as the acrylic hydrophilicitizing agent in each one of the Examples disclosed in the publication, although the use of equivalent amounts of methacrylic acid, acrylamide and other acrylic hydrophilicitizing agents is anticipated.
The publication of Choi (2004) discloses the use of "ethylenically unsaturated hydrophilicitizing agents" in graft polymerisation. These other hydrophilicitizing agents are disclosed as hydrophilicitizing agents that are polymerisable by addition polymerisation to a thermoplastic polymer and are hydrophilic as a consequence of containing carboxyl (—COOH), hydroxyl (—OH), sulfonyl (SO3), sulfonic acid (—SO3H) or carbonyl (—CO) groups. No experimental results concerning the chemical and physical properties of graft polymers prepared by a method using these other hydrophilicitizing agents is disclosed.
The publication of Choi (2005) discloses a non-woven sheet of polyolefin fibres where opposed surfaces of the sheet are hydrophilic as a consequence of an acrylic graft polymerisation. The properties of the sheet are asymmetric, the ion exchange coefficient of the two surfaces being different. The method used to prepare these asymmetric acrylic graft polymerised non-woven polyolefin sheets comprises the steps of immersing the substrates in a solution of benzophenone (a photoinitiator), drying and then immersing the substrate in a solution of acrylic acid prior to subjecting to ultraviolet (UV) irradiation. The irradiation may be performed when the surfaces are in contact with either air or an inert atmosphere.
The publication of Gao et al (2013) discloses a method of preparing a radiation cross-linked lithium-ion battery separator. In an example, a porous polyethylene membrane is immersed in a solution of benzophenone and triallyl cyanurate in dichloromethane. The immersed membrane is dried at room temperature before being immersed in a water bath at 30° C. and irradiated on both sides using a high-pressure mercury lamp for three minutes.
The objective of the majority of these prior art methods is to improve the adhesion, biocompatibility, printability or wettability of the surface of a substrate. These improvements to surface characteristics are to be distinguished from the use of UV-initiated polymerisation to modify the permeability of preformed microporous polyolefin substrates, such as the substrates described in the publications of Fisher et al (1991) and Gillberg-LaForce (1994).
It is well-known to prepare thin film composite membranes to modify the permeability of a preformed microporous polyolefin substrate. For example, the publication of Jones (1990) discloses a composite permselective membrane comprising an ultrathin semipermeable layer comprising a polybenzimidazole polymer in occluding contact with at least one surface of a microporous polymer support layer. The membranes are asserted to provide better combinations of flux and rejection rates in reverse osmosis processes than do conventional semipermeable membranes of polybenzimidazole polymer alone.
The publication of Callahan and Johnson (1990) discloses a composite membrane having a microporous support which is coated with a UV curable polymer composition having a sufficiently high viscosity to prevent pore filling upon coating and curing.
The publication of Gillberg-LaForce and Gabriel (1991) discloses a pore modified microporous membrane which is made by a process of incorporating a polymerizable vinyl hydrophilicitizing agent within the pores of a microporous membrane followed by polymerization to secure the resulting polymer within the pores. The process is stated to be particularly suitable for modifying a hydrophobic microporous membrane with a hydrophilic polymer, as occurs for example when polyacrylic acid is secured into the pores of a polypropylene microporous membrane.
The publication of Callahan and Johnson (1992) discloses a composite membrane having a microporous support which is coated with a UV curable polymer composition having a sufficiently high viscosity to prevent pore filling upon coating and curing.
The publication of Cussler et al (1992) discloses a process for modifying the properties of a hydrophobic microporous membrane which includes the steps of treating a hydrophobic microporous membrane with a surfactant to render the membrane hydrophilic, permeating the membrane with a polyol, and crosslinking the polyol to yield a hydrophilic microporous membrane having pores filled with an aqueous gel. The modified membranes are asserted to be useful in carrying out chromatographic separations.
The publication of Donato and Phillips (1993) discloses a composite membrane having a microporous support which is coated with a polymer selected from the group consisting of polyethylene oxide, polyacrylic acid, poly(methyl methacrylate) and polyacrylamide wherein there is no pore filing of the microporous support. The publication of Donato (1994) discloses a composite membrane having a microporous support coated with an aqueous polyeurethane dispersion composition. The publication of Donato and Phillips (1994) discloses a composite membrane having a microporous support which is coated with a polymer composition and a contact adhesive layer applied to said polymer.
For the most part, the methods of preparing composite membranes disclosed in these publications use UV initiated polymerisation to form polymers in situ. Methods of adhering dissimilar preformed polymers to the surface of the microporous polyolefin substrates are less well known.
It is an object of the present invention to provide a method of decreasing the hydrophobicity of preformed microporous poly(ethylene) sheets and thereby provided modified microporous polyethylene sheets suitable for use in the preparation of water permeable asymmetric composite membranes. It is an object of the present invention to provide an asymmetric composite membrane suitable for use in the recovery or removal of water from dairy and other feed streams. These objects are to be read in the alternative with the object at least to provide a useful choice.

Asymmetric composite membranes and modified substrates used in their preparation


A support structure for supporting a pressurizer safety and relief valve system above a pressurizer vessel comprises a ring girder which surrounds the pressurizer vessel wall and is supported on pairs of lugs spaced about the pressurizer wall. The ring girder supports columns for the system at locations spaced from the lugs, and support means in the form of frames distribute the load on the lugs. The frames have upper and lower bars and outer vertical bars, with spaced inner vertical bars. Biasing means are provided to attach the frames to the lugs and distribute the forces, and a projection on the ring girder fits into a groove of a support bar on the upper bar of the frame.

The present invention relates to a support structure for supporting a pressurizer safety and relief valve system above the pressurizer vessel of a pressurized water nuclear reactor power plant.
A pressurizer vessel is used in a pressurized water nuclear reactor power plant to provide relief for coolant overpressure. Usually, such a pressurizer vessel is a vertical, thick walled vessel having a hemispherical top, or dome, and having its bottom end supported by a flanged skirt.
Relief valves are provided for a first level of overpressure protection and safety valves are provided for a second level of overpressure protection. The unit thus comprises a pressurizer relief and safety valve system that has a manifold or header assembly and connections for nozzles in the head of the pressurizer, with piping and discharge outlets provided.
A pressurizer safety and relief valve system which is supported above the pressurizer on columnar supports is described in co-pending application Ser. No. 229,363, now U.S. Pat. No. 4,426,350, dated Jan. 17, 1984, entitled "Improved Valve Support Arrangement for Pressurizer in a Nuclear Power Plant", filed Jan. 29, 1981, in the names of Martin J. Zegar, et al., assigned to the assignee of the present invention, which application is incorporated by reference herein. In said co-pending application, an arrangement for supporting the pressurizer safety and relief valve system is described wherein a common header is secured to the side wall of the pressurizer vessel by the use of columnar supports that are secured to the header and are secured to the pressurizer side wall by the use of lug means. The lug means which are attached to the pressurizer side wall must be capable of supporting the load of the columnar supports and pressurizer safety and relief valve system and could become overstressed for some combinations of loads, and also must be of sufficient dimensions if they are to support the necessary loads.
In the system described in said co-pending application, the main support structure comprises four columnar members equally spaced around the pressurizer vessel and coupled to the lugs on the vessel side wall. A fixed type connection between the lugs and the columnar supports can cause undesirable concentrations of stress and difficulties in fabrication and construction.
It is an object of the present invention to provide a support structure for a pressurizer safety and relief valve system that will ease the stress strain condition of the lugs supporting the system.
It is another object of the present invention to provide a frame for a ring girder, for supporting a pressurizer safety and relief valve system that is sufficiently rigid so as not to bend along a horizontal axis to provide a uniform load on the upper support area of a lug, and sufficiently flexible to provide for possible deflection due to forces caused by temperature variations.

Stirrup-type support structure for nuclear power plant pressurizer valves


A coin processing apparatus which comprises a coin selecting section for determining whether put-in coins are genuine or not and denominations of the coins determined as genuine, a coin accommodating section for accommodating the genuine coins selected by the coin selecting section, and a coin payment section for paying out coins having denominations corresponding to an amount of change from the coin accommodating section, in which the coin accommodating section is formed at the lower portion of a single housing and the coin selecting section is formed by an inner wall face of the housing having a part of the coin selecting section and by a cover having another part of the coin selecting section on its back face, the cover being adapted for covering the inner wall face of the housing, and in which a rib is formed on an upper surface of the housing, for guiding liquid fell on the upper surface of the housing to both sides of the housing and, the coin selecting section is provided with a liquid collector, a discharge path and a guide groove for rapidly discharging liquid poured into a coin slot to the outside of the coin processing apparatus.

1. Field of the Invention
The present invention relates to coin processing apparatus used in vending machines, money exchangers, service devices, etc., and adapted to separate and accommodate put-in coins in accordance with denominations and pay out those separated and accommodated coins as change.
2. Description of the Related Art
Generally, vending machines, money exchangers, service devices, etc., include a coin processing apparatus for separating and accommodating put-in coins in accordance with denominations and paying out those coins as change.
FIG. 17 is a decomposed perspective view of a conventional coin processing apparatus 100.
This coin processing apparatus 100 is mainly composed of three independent units, i.e., a coin payment unit 102 configured by a housing having a substantially -shaped sectional configuration and provided with a pay-out mechanism 101 at its lower portion for paying out change, a coin accommodating unit 103 composed of coin tubes to be mounted above the pay-out mechanism 101 of the coin payment unit 102 and, a coin selecting unit 104 to be mounted above the coin accommodating unit 103.
Among these units constituting the coin processing apparatus 100, the coin accommodating unit 103 and the coin selecting unit 104 are arranged to be mounted within the coin payment unit 102 as shown by arrows when they are used. In other words, the coin accommodating unit 103 and the coin selecting unit 104 are formed separately from the coin payment unit 102 and are configured to be detachably mounted to the coin payment unit 102.
When a coin is put into the coin processing apparatus 100, the coin first enters into the coin selecting unit 104 and checked whether a genuine one or not, and if it is a genuine one, it is classified in accordance with its denomination. The genuine coin is then accommodated in accordance with its denomination in the coin accommodating unit 103 disposed below the coin selecting unit 104. When denominations of change are determined, coins accommodated in the coin accommodating unit 103 are paid out for an amount corresponding to the sum of change from the lower portion of the coin processing apparatus 100 by the coin payment mechanism 101 formed at the lowermost position of the coin payment unit 102.
Since the conventional coin processing apparatus 100 is configured such that each of the units are separately provided and the coin selecting unit 104 having a complicated configuration is detachably mounted to the coin payment unit 102, the entire configuration of the coin processing apparatus 100 becomes extremely complicated. Further, a large number of mechanical and electrical parts (connectors, connecting elements, etc.) for connecting the respective units (particularly the coin selecting unit) to the coin payment unit 102 are required and thus the manufacturing cost of the entire coin processing apparatus 100 becomes high.
According to the conventional coin processing apparatus 100 as described above, complicated coin paths for selection of coins are formed in the coin selecting unit 104. Since coins roll down along the coin paths, if liquid such as soft drink or cleaning detergent is poured into a coin slot of a vending machine, for example, the liquid will flow along the upper surface of the coin processing apparatus 100 and along the coin paths of the coin selecting unit 104. When adhesive liquid such as cleaning detergent flows through the coin paths in the coin selecting unit, it will adhere to selection units such as coin separating levers disposed in the coin paths whereby the normal operation of such selection units is hindered and hence the selecting function of the coin processing apparatus 100 is deteriorated. When such adhesive liquid adhering to the selection units is dried, the operation of these selection units may be locked and the function of the coin processing apparatus 100 itself may be stopped. Further the liquid may flow into the coin tubes which function as the coin accommodating unit 103 and also into the coin payment unit 102 causing short-circuit of electronic circuits and a payment motor which results in a malfunction of the apparatus.


For the setting of yarn clearers for fineness control the number of permissible false alarms for a given yarn length, i.e. the permissible alarm frequency, is fixed. During the clearing process measured values of the fineness are recorded continuously, and their distribution is determined. From this distribution of the measured values and from the predetermined permissible alarm frequency, response limits are fixed independently on the basis of statistical regularities. Thus, the yarn clearers can be set automatically in an optimum manner on the basis of actual production circumstances, and mistakes between yarns differing only slightly in their fineness cannot occur.

The invention relates to a method for setting the sensitivity limits of electronic yarn clearers taking into account the fineness or count of the yarn to be cleared, the fineness being measured during the clearing process, and the result of this measurement being compared with set tolerance limits.
It is known (see, for example, USTER News Bulletin No. 29/August 1981 "Das Uster-System der Garnfehlerkontrolle" (The USTER System of Yarn Fault Control, USTER being a registered trademark of Zellweger Uster AG), the disclosure of which is hereby incorporated by reference in its entirety), that an essential precondition for accurate and reproducible yarn clearing is the establishment of a basic setting for the yarn batch to be cleared. As described in the USTER News Bulletin, yarn clearing ensures that yarn quality is suitable for subsequent processing by deleting and eliminating yarn faults. For example, typical yarn faults include deviations from a normal yarn cross-section which result in the occurrence of, for example, short thick places (e.g., a fault length of yarn less than 8 cm having a cross-section approximately twice that of the yarn), long thick places (e.g., a fault length greater than 40 cm having a cross-section approximately twice that of yarn), or thin places (e.g., having a variable fault length and a cross-section approximately 30 to 70% that of the yarn).
Although these yarn faults seldom occur, they can significantly affect the appearance of a fabric woven or processed from the yarn and/or negatively influence subsequent processing of the yarn. However, the number of faults which are actually eliminated over a given length of yarn is often limited to avoid subsequent processing difficulties. For example, as described in the aforementioned 1981 USTER News Bulletin, thin place faults thinner than 60% of normal yarn cross-section can be extracted when they are longer than approximately 40 cm, with faults shorter than 40 cm in length only being extracted if the cross-section is proportionately less than 60% of normal cross-section. Extraction of a maximum of 1-2 thin places per 100,000 cm of yarn is, for example, preferred for cotton spun yarns. Different settings can be provided for other faults, such as thick place faults and for different materials. Thus, the most important value which must be taken into account in establishing a basic setting for a yarn batch being cleared is the yarn count or fineness, which is used to fix the response limits (=tolerance values) which if exceeded give rise to a clearer step being triggered.
Since yarns of different fineness are normally spun in spinning mills, mistakes can occur, in particular in cases where the yarns differ only slightly in fineness from one another. However, even slight differences in the fineness lead to visible stripes in the woven or knitted fabric, making the product unusable. Yarn monitoring units, which measure the yarn fineness and trigger off an alarm or stop the production if the set tolerance limit is exceeded are available in order avoid this risk of confusion. However, every yarn naturally exhibits a certain unevenness, and since certain fineness deviations cannot be entirely avoided in the normal production process, difficulties are encountered in the setting of the response limits. For example, if the tolerance limits are set too narrow, false alarms frequently occur. On the other hand, if the tolerance limits are set too wide, not all faults which may be detected and eliminated without exceeding the number of correctable faults per unit length of yarn can be detected.

Method for setting the sensitivity limits of electronic yarn clearers, and device for carrying out the method


A mold element comprised of metallic pellets and a bonding agent to bond together the metallic pellets. The metallic pellets can be a mixture of ferrous shot and ferrous grit and the bonding agent can be a solvent or a resin. An associated method is also provided.

1. Field of the Invention
The present invention relates to a mold element having a plurality of metallic elements secured to each other by means of a bonding agent, and a related method.
2. Brief Description of the Prior Art
It is well known to provide molds for casting certain objects. These molds have a cavity which is formed therein to correspond substantially to the external shape of the casting to be cast. The cavity is adapted to be filled with casting material through a suitable pouring hole.
Many times there is a need to form internal passages in metal castings. In order to accomplish this a "core" or "cores" are placed in the mold. Mold cores are traditionally produced from loose, divided grains which are "bound" or "adhered" together in a manner such that a solid structure is fashioned. These materials afford dimensional and structural stability and resist the tendency to combust. See U.S. Pat. Nos. 3,008,200 and 3,228,074. Once the molten metal is poured around such a core structure and the metal again solidifies, the core must be removed in such a fashion that the internal passage is retained within the cast shape. The most wide-spread method of core removal is disintegration of the core by "knock-out" methods using vibration and/or shock and by impact blast cleaning through the use of a focused media stream.
Traditionally, mold and mold core materials have consisted of various silica grains, ceramic grains, clays, and the like, held in place by "binders" of numerous compositions, such as a phenolic resin. See U.S. Pat. No. 3,228,074. It has also been known to magnetize the core material See U.S. Pat. No. 3,008,200.
Use of metal-metal oxide compositions for various articles is known. See, for example, U.S. Pat. No. 4,255,193 and United Kingdom Patent No. 321,394.
There is a present demand for greater casting integrity and for improved casting density in thin cast sections. There remains therefore, a very real and substantial need for inexpensive molds and mold cores that may be provided in a range of sizes that retains a high degree of structural accuracy.


A computer-implemented method includes receiving a first 3D model and a second 3D model, wherein the first 3D model represents a first set of objects and the second 3D model represents includes a second set of objects. The computer-implemented method further includes scanning each of the first 3D model and the second 3D model to identify the first set of objects and the second set of objects, wherein the first set of objects and the second set of objects have at least one common object. The computer-implemented method further includes comparing the first set of objects to the second set of objects to yield one or more differences. The computer-implemented method further includes sorting each of the one or more differences based on a set of rules to yield a list of differences. A corresponding computer system and computer program product are also disclosed.

The present disclosure relates generally to 3D models and in particular to detecting and presenting differences between 3D models.
The development of computer vision has led to an increase in the ability to recover the three-dimensional shape and appearance of objects in imagery. Currently, there exist various domains for which 3D environmental comparisons are utilized. For example, 3D models may be used for post disaster assessment of infrastructure systems and buildings by comparing a 3D model of pre-event conditions to a 3D model of post-event conditions. However, the number of differences that may be detected between 3D models can be excessive and overwhelmingly time consuming to review.

Detection and presentation of differences between 3D models

International Business Machines Corporation


The present invention inter alia provides a method, and use thereof, of predicting CV events such as AMI, ACS, stroke, and CV death by determining the concentration of at least one ceramide of Formula I or one lysophospholipid of Formula II and/or III and at least one lysophospholipid of Formula IV, V, VI, VII and/or VIII in a biological sample and comparing those concentrations to a control. Finding an increased concentration of the at least one Formula I ceramide or Formula II and/or III lysophospholipid and a decreased concentration of the at least one Formula IV, V, VI, VII and/or VIII lysophospholipid indicates that the subject has an increased risk of developing one or more CV events. The present disclosure also provides a method, and use thereof, of diagnosing subjects suffering acute ischemia. Also provided are kits and compositions comprising the same for use in predicting and/or diagnosing CV events.

Current CV markers, LDL-C, LDL-P, and total cholesterol, are not able to identify coronary artery disease (CAD) patients or subjects that have an elevated risk for CV events, such as AMI, ACS, stroke and CV death, from patients having more stable disease.
The term myocardial infarction pathologically denotes the death of cardiac myocytes due to extended ischemia, which may be caused by an increase in perfusion demand or a decrease in blood flow. The event is called "acute" if it is sudden and serious. Diagnosis of AMI is determined by a high clinical suspicion from history and physical examination, in addition to changes in cardiac biomarkers (creatinine kinase MB [CK-MB], troponins, and myoglobin) and electrocardiogram (ECG) findings. Imaging techniques, such as two-dimensional echocardiography, are also useful in demonstrating myocardial dysfunction. Current CV markers, LDL-C, LDL-P, and total cholesterol, fail to predict the likelihood of the cardiovascular disease (CVD) event occurring.
Acute coronary syndrome (ACS) is a term used for any condition brought on by sudden, reduced blood flow to the heart. The first sign of acute coronary syndrome can be sudden stopping of the heart called cardiac arrest. Acute coronary syndrome is often diagnosed in an emergency room or hospital with the same cardiac biomarkers or electrocardiogram (ECG) as AMI, that provide evidence on damaged heart tissue or problems in the heart's electric activity.
A stroke is the loss of brain function due to a disturbance in the blood supply to the brain, depriving brain tissue of oxygen and food. Within minutes, brain cells begin to die. A stroke may be caused by a blocked artery (ischemic stroke) or a leaking or burst blood vessel (hemorrhagic stroke). Some people may experience a temporary disruption of blood flow through their brain (transient ischemic attack, or TIA). Strokes are usually diagnosed by brain imaging and carrying out physical tests.
Sudden cardiac death (SCD) is a sudden, unexpected death caused by loss of heart function, also named as sudden cardiac arrest (SCA). Sudden cardiac arrest is not a heart attack (myocardial infarction). Heart attacks occur when there is a blockage in one or more of the coronary arteries, preventing the heart from receiving enough oxygen-rich blood. In contrast, sudden cardiac arrest occurs when the electrical system to the heart malfunctions and suddenly becomes very irregular. The heart beats dangerously fast. Ventricular fibrillation may occur, and blood is not delivered to the body. In the first few minutes, the greatest concern is that blood flow to the brain will be reduced so drastically that a person will lose consciousness. Death follows unless emergency treatment is begun immediately. Sudden cardiac arrest happens without warning and is rarely diagnosed with medical tests as its happening. Instead, SCA often is diagnosed after it happens, by ruling out other causes of a person's sudden collapse.
AMI, ACS, stroke and sudden cardiac death are diagnosed in acute stage, but predictive markers are not available. The risk factors behind these events are for example, age, hypolipidemia, hypertension, smoking, diabetes, CAD or previous heart attack. Yet, no diagnostic test that could predict the events exists, and cardiovascular diseases are the leading cause of death worldwide. Furthermore, the costs to society that are associated with CVD are higher than for any other group of diseases. The same tests that are used for diagnosing CVD are utilized in predicting the events. Today the most innovative approach is to use LDL-C, LDL-P, HDL-C, Lp(a), Lp-PLA2 or CRP. However, none of the listed lipid based markers (LDL-C, LDL-P, Lp(a), Lp-PLA2) provide clinically useful predictive information allowing stratification to aid physicians. CRP has been promising in the research setting, however it has proven to be unspecific (CRP is an acute phase reactant that can react to many different stimuli leading to highly variable test results) and thus CRP values are difficult to interpret in the clinical use. There is an unmet need for a diagnostic test that could predict CV events, such as AMI (acute myocardial infarction), ACS (acute coronary syndrome), stroke and CV death.
Ischemia is a restriction in blood supply to tissues. When blood flow to the heart is reduced, the heart does not receive enough oxygen and nutrients, and the condition is called myocardial ischemia or cardiac ischemia. Myocardial ischemia can damage the heart muscle leading, e.g., to a heart attack. Therefore, an early detection of ischemia is crucial for survival of patients. There is an unmet need for diagnostic markers that could identify patients with early stage acute ischemia that can be either symptomatic or asymptomatic (silent ischemia).
The ceramide and lysophospholipid (LPL) based risk stratification offers improved performance compared to any other lipid based biomarker today.
Alterations in specific LPLs, i.e. lysophosphatidylcholines (LPCs), have been observed in primary prevention settings related to CAD (BG Medicine disclosures EP2293077 and U.S. Pat. No. 8,321,154 as well as Fernandez et al., Plos one, 2013, 8(8):e71846). However, these studies have not provided information regarding the outcome of the patients, and therefore do not report those markers that could be used to identify CAD patients at high risk for future outcome events. Moreover, these studies did not investigate the association of ceramides related to CAD and did not investigate LPLs other than LPCs.
A large group of lipid molecules, including certain ceramides and LPCs, and ratios calculated from two lipid molecules have been identified for predicting CV outcomes in CAD patients who are undergoing statin treatment or who are not undergoing statin treatment or for identifying high-risk CAD patients or predicting whether a subject is at risk for developing CV events. However, there remains a desire for improved methods of predicting the risk of a patient developing a CV event, such as AMI, ACS, stroke, and CV death.
The current application recognizes value of both ceramides (Tarasov et al: J Clin Endocrinol Metab., 2014, 99(1):E45-52) and lysophospholipids (Goncalves et al., Arterioscler Thromb Vasc Biol. 2012 June; 32(6):1505-12 PMID:22499993) in development of CVD. It has earlier been reported that blood ceramide levels can be lowered by certain lipid lowering compounds including statin drugs (Tarasov et al., J Clin Endocrinol Metab., 2014, 99(1):E45-52) and newly developed Lp-PLA2 inhibitors may potentially affect Lp-PLA2 mediated hydrolysis of low-density lipoprotein-oxidized phospholipids leading to generation of LPLs. Therefore, Cer/LPL-ratio can provide significant clinical improvement compared for instance to certain previously disclosed Cer/Cer-ratios as Cer/LPL-ratios can offer a better monitoring opportunity for the drug treatment effects; Cer can be affected by statin type of cholesterol lowering compounds, while LPL component may be controlled by drugs affecting for instance Lp-PLA2 activity.
This application provides new ceramide and LPL markers with superior AUC, sensitivity and specificity. This superior performance could not have been predicted from art known biomarkers.

Use of ceramides and LPLs in diagnosing CVD


An electronic device includes a display module having a hole and including a plurality of first light sources disposed along a perimeter of one side of the display module and a plurality of second light sources which are positioned opposite the plurality of first light sources and are disposed along a perimeter of the other side of the display module, a shaft rotatably inserted into the hole of the display module, a movement which is positioned in the rear of the display module, is connected to one end of the shaft, and provides a rotational force for the shaft, a window which is separated from the display module and is positioned in front of the display module, and a hand which is positioned between the display module and the window and is fixed to the other end of the shaft.

Field of the Invention
The present disclosure relates to an electronic device, and more particularly, to an electronic device involved in a wearable watch, in which analog and digital are combined.
Discussion of the Related Art
Terminals may be generally classified into mobile/portable terminals and stationary terminals based on a mobility. The mobile terminals may also be classified into handheld terminals and vehicle mounted terminals depending on whether or not a user can directly carry the terminal.
Mobile terminals have increasingly more functions. Examples of the functions include data and voice communications, capturing images and video using a camera, recording audio, playing music files using a speaker system, and displaying images and video on a display. Some mobile terminals include additional functionality which supports game playing, while other terminals are configured as multimedia players. More recently, the mobile terminals have been configured to receive broadcast and multicast signals which permit viewing of content such as videos and television programs.
As the mobile terminals have increasingly more functions, the mobile terminals have been implemented as multimedia players of multiple functions having capturing images and video, playing music files or video, game playing, receiving broadcast, etc.
Efforts are ongoing to support and increase the functionality of mobile terminals. Such efforts include software and hardware improvements, as well as changes and improvements in the structural components.
A study on wearable electronic devices the user directly wears is being recently carried out. For example, an attempt has been made to study wearable electronic devices of a glass type, a watch type, etc.
Because the wearable electronic device has to arrange necessary electronic components in a limited space while satisfying design requirements, a need for an optimum design of the wearable electronic device is increasing.

This patent is subject to a terminal disclaimer.


A transmission and reception system includes a transmitting unit and a receiving unit. The transmitting unit transmits a transmission signal including at least a data signal and a clock signal. The data signal is transmitted at a random timing. The receiving unit includes a delay unit for delaying the data signal by a predetermined time period in synchronous with the clock signal. A buffer unit operates in synchronous with the clock signal, temporally stores the data of the delayed data signal and outputs the data. A control units operates in synchronous with the clock signal, controls the storing and outputting operations of the buffer unit. A reset unit operates in synchronous with the clock signal, supplies a reset signal at start timing of the data signal to reset the control unit.

1. Field of the Invention
The present invention relates to a receiving circuit and a transmission/reception system utilizing the same, more particularly, to a receiving circuit reset upon reception of burst data in a transmission/reception system in which the burst data and a clock signal are transmitted at least from a transmitting unit to a receiving unit.
2. Description of Related Art
FIG. 1 is a block diagram showing a conventional transmission/reception system. The system includes a transmitting unit 102 and a receiving unit 104 containing a receiving circuit 106. The transmitting unit 102 transmits a burst data signal DATA0 and a clock signal CLK to the receiving unit 104. The receiving circuit 106 receives these signals and outputs a data signal DATA1 from these signals. The system is used in a computer system, for instance. In this case, a computer corresponds to the transmitting unit and a unit for receiving data from the computer corresponds to the receiving unit.
FIG. 2 is a block diagram showing the receiving circuit 106. The receiving circuit 106 includes a detector 114, a control unit 118, a buffer unit 122, and an output unit 124. All of these units are supplied with the clock signal CLK and operates in synchronously with the clock signal. The detector 114 receives the data DATA0 from the transmitting unit 102 to detect whether there is the burst data, and generates a signal DP indicating the presence of burst data when it is determined that the burst data is present. The buffer unit 122 also receives the data DATA0 from the transmitting unit 102 in response to a signal DRS from the control unit 118 to store it therein. Also, the buffer unit 122 generates a signal SD while the burst data is stored or held therein. The output unit 124 outputs the stored data from the buffer unit 122 as data DATA1 in response to a signal DOE from the control unit 118. The control unit 118 controls the detector 114, the buffer unit 122, and the output unit 124. More particularly, the control unit 118 receives the signal DP from the detector to generate the signal DRS indicating that the burst data is to be stored in the buffer unit 122. The control unit 118 generates the signal DOE in response to the signal SD from the buffer unit 122 and the signal DP from the detector 114.
FIG. 3 is a block diagram of the control unit 118. The control unit 118 includes a control logic unit 118-1 and a status register 118-2. The status register 118-2 stores the status of the receiving circuit 106 in response to the clock signal CLK. The control logic unit 118-1 generates the control signals DRS and DOE in accordance with the signals DP and SD and the current status stored in the register 118-2.
FIG. 4 is a state transition diagram of the receiving circuit 106. The receiving circuit 106 waits for burst data in the state S2. When the burst data is transmitted from the transmitting unit 102 to the receiving unit 104, the state changes to an on-reception state S4 because the detector detects the burst data to issue the signal DP to the control unit 118. When the transmission of burst data is completed, i.e., when the signal DP indicates that there is no data, the control unit 118 changes the state to the data output state S6 in accordance with the signal SD to issue the signal DOE to the output unit 124 so that the data stored in the buffer unit is outputted as the DATA1.
However, in this conventional receiving circuit, in a case where the clock signal is distorted due to the application of noise on the transmission path, a clock signal having a high clock rate would be determined to be inputted. When there is inputted to the receiving circuit such a clock signal having a shorter clock pulse than a width of the clock pulse contained in the normal clock signal, the control circuit 118 transits to an abnormal state so that the operation is in a fault state. In this case, flip-flop circuits needs to be reset, and unless the flip-flop circuit are not reset the control unit is not recovered to the normal state.

Receiving circuit reset upon reception of burst data and transmission/reception system utilizing the same


This invention relates to methods, a computer program product and apparatuses in the context of frame buffering. A buffering time for one or more frames received by a frame buffer is determined based on a specific buffering time associated with a specific frame and on information representative of a specific amount of data stored in the frame buffer.

Frame buffers are inter alia deployed in the context of video and audio processing, such as for instance in Voice over IP (VoIP) systems.
In general, network jitter and packet loss conditions can cause degradation in quality for example in conversational speech services in packet switched networks, such as the Internet. The nature of the packet switched communications can introduce variations in transmission of times of the packets (containing frames), known as jitter, which can be seen by the receiver as packets arriving at irregular intervals. However, an audio playback device requires constant input to maintain good audio quality, and no interruptions can be allowed. Thus, if some packets/frames arrive at the receiver after they are required for playback, the decoder may have to consider those frames as lost and perform error concealment.
Typically, a fixed buffer to manage network jitter can be utilised to store incoming frames for a predetermined amount of time (specified e.g. upon reception of the first packet of a stream) to hide the irregular arrival times and provide constant input to the decoder and playback components.
This approach is illustrated in FIG. 1, which shows timelines of a transmitter TX, a receiver RX and a decoder DEC, respectively. On the uppermost timeline, the regularly spaced transmission points of time of frames are illustrated by vertical lines. Similarly, in the middle timeline, the reception points of time of the transmitted frames are illustrated by vertical lines, wherein the reception points of time, for instance due to network jitter, are irregularly spaced. Finally, the lowermost timeline shows the decoding points of time of the frames, which are, due to the fact that the decoder DEC needs to provide constant input to the playback device, regularly spaced. Therein, the relationship between transmission, reception and decoding points of time of the same frame is illustrated by means of dashed arrows. As can be seen by comparing the reception point of time of the first received frame in the middle timeline with its associated decoding point of time in the lowermost timeline, the jitter buffer introduces an additional delay increasing the end-to-end delay since the received frames are stored in the jitter buffer before the decoding process.
A shortcoming of this basic approach is that a jitter buffer management scheme requiring fixed playback timing is inevitably a compromise between low enough buffering delay and low enough number of delayed frames, and finding an optimal trade-off can be a difficult task. For example, in the situation shown in FIG. 1, the frame sent at 100 ms arrives at the receiver RX after it is needed for further processing by decoder DEC, and therefore it needs to be replaced with error concealment (i.e. a "late loss" occurs, shown as a dashed vertical line).
There may exist special environments and applications where the amount of expected jitter can be estimated to remain between predetermined limits. In general, however, the network delay associated with jitter can vary from a scintilla of time to hundreds of milliseconds within the same session. Using a jitter buffer management scheme requiring fixed playback timing with an initial buffering delay set to a large enough value to cover the jitter, according to an expected worst case scenario, would keep the number of delayed packets in control. However, at the same time there may arise a risk of introducing an end-to-end delay which can be too long to enable a natural conversation. In this case applying a fixed jitter buffer management scheme may not be a practical choice in most audio transmission applications operating over a packet switched network, e.g. in VoIP over the 3GPP IP Multimedia Subsystem (IMS).
In contrast to a fixed jitter buffer management scheme, an adaptive jitter buffer management scheme can be used to dynamically control the balance between short enough delay and low enough number of delayed frames. In this approach, an entity controlling the jitter buffer constantly monitors the incoming packet stream and adjusts the buffering delay (or buffering time, these terms are used interchangeably) according to observed changes in the network delay behaviour. If the transmission delay seems to increase or the jitter becomes worse, the buffering delay may need to be increased to meet the network conditions. In the opposite situation, where the transmission delay seems to decrease, the buffering delay can be reduced, and hence, the overall end-to-end delay can be minimised.
One of the challenges in adaptive jitter buffer management is reliable estimation—or actually prediction—of the transmission characteristics. Although adaptation based on the reception statistics of most recent packets usually gives a reasonable estimate on the short-term network behaviour, it may be impossible to avoid the fact that some frames arrive after their scheduled decoding time—i.e. too late for normal decoding, especially when applying relatively strict buffering delay requirement.
Jitter buffer adaptation during active speech requires additional processing to shorten or extend the speech signal (i.e. time scaling, also known as time warping) to maintain good voice quality and intelligibility. For example, suitable methods are disclosed in documents WO 03/021830 and U.S. 2006/0056383. To avoid complex time scaling, a commonly used method for jitter buffer management is to perform the adaptation during comfort noise signal periods typically at beginning of a new talk spurt (i.e. at a speech onset). This approach can be expected to provide low complexity adaptation functionality with high quality, since the comfort noise signal does not carry information that is important for intelligibility or actual voice quality. The minor drawback of the onset adaptive jitter buffer management is that even though the network analyser detects changes in the delay characteristics, the jitter buffer adaptation needs to wait for the next speech onset to take place. However, jitter buffer management solutions apply the onset adaptive approach as part of the adaptation functionality. Where the basic approach is to re-estimate the required buffering time and perform adaptation at each speech onset, while only urgent adaptation steps are taken during active speech.
A basic (adaptive) jitter buffer management approach uses the statistics on the current number of frames in the buffer as an indication of the buffer status. If, for example, the number of frames in the buffer falls below a predetermined (or adaptively determined) threshold, an adaptation step to increase the buffering time can be initiated to take place at the next speech onset in order to decrease the risk of subsequent frames arriving too late. If, however, the number of frames in the buffer grows above another predetermined (or adaptively determined) threshold, an adaptation step to decrease the buffering time can be initiated to reduce delay for improved interactivity.
An alternative approach may use statistics computed over several frames instead of considering only single frames. For instance, the number of instances when the number of frames in the buffer falls below or above predetermined limits over an analysis window consisting of several frames (or time corresponding to several frames' duration) may be counted. Equally well, the average number of frames in the buffer may be considered over an analysis window as the indication of buffer status which is used for controlling the jitter buffer management operation.
When considering the buffer status indication based on the number of frames, one possible way to apply adaptation is to estimate the target buffer level as a number of frames, and then to wait for the selected number of frames to arrive and accumulate in the buffer before decoding (or playing back) the first frame of a talk spurt at the point a speech onset is detected.
This is illustrated in FIG. 2, wherein the structure of FIG. 2. In this example, the speech onset frame is the one transmitted at 20 ms (and received approximately at 36 ms), and the decoding & playback is started after three frames (the first one to be decoded plus two other ones providing assumed 40 ms of jitter protection) have arrived at the buffer. In this example this approach would result in approximately 52 ms buffering delay for the first frame of the talk spurt.
A considerably different jitter buffer management approach is to indicate the buffer status based on the buffering time. The buffering time may for instance be an observed buffering time recorded when a frame is passed to the decoder (or playback device), or a predicted buffering time computed when a frame is received. Methods for predicting a buffering time are for instance disclosed in document WO 2006/044696. As in the aforementioned approaches, which are based on the number of frames stored in the buffer, the buffer adaptation approach based on the buffering time may also consider statistics computed over several frames instead of keeping track of single frames only.
Jitter buffer management based on the buffering time may for instance be performed by estimating the required buffering time (e.g. in milliseconds), and applying this buffering time for the first frame of a talk spurt when a speech onset is encountered.
The example in FIG. 3 illustrates this approach. In this case, the estimated buffering delay of 40 ms is applied for the first frame of the talk spurt simply by directly delaying its decoding & playback by 40 ms.
The previously described jitter buffer management approaches for speech onset handling can be expected to provide, on average, approximately equal performance (both in terms of delay and late-loss). However, there can be special cases where these approaches may fail, either by introducing unnecessarily high buffering delay until the next adaptation step, or by providing inadequate jitter protection, leading to an unnecessarily high rate of late-loss frames.
For instance, problematic cases for jitter buffer management based on the number of frames in the buffer are scenarios where the frame triggering the start of the decoding & playback is an "outlier"—i.e. it arrives too early or too late (with respect to the subsequent frames). In the former case, which is illustrated in FIG. 4, buffer adaptation may provide too short buffering time causing high late-loss rate, while in the latter case, which is illustrated in FIG. 5, the buffering time would be unnecessarily high (and there may be also a risk of buffer overflow) until the next adaptation point.
If the "trigger frame", denoting the frame triggering the start of the decoding and playback, is lost in the transmission path, the jitter buffer manager will wait until the next frame is successfully received before starting the decoding & playback. In case of a long lost burst this may obviously increase buffering time significantly.
On the other hand, for jitter buffer management based on the buffering time, the critical frame is the actual speech onset frame: if it arrives "early", the buffering time may be too short (see FIG. 6), causing increased late-loss rate. If it arrives "late", the buffering time may become unnecessarily long (see FIG. 7).


The present invention provides a measuring apparatus comprising: a light source unit; and a sensor unit configured to sense a plurality of properties to be measured with respect to an object by using light emitted from the light source unit, wherein the sensor unit comprises: a first region configured to sense a first property to be measured among the plurality of properties to be measured by sensing the light within a first wavelength range; and a second region configured to sense a second property to be measured among the plurality of properties to be measured by sensing the light within a second wavelength range that is different from the first wavelength range.

X-ray radiography, ultrasonic imaging, computed tomography (CT) and the like, which are diagnostic imaging used in a clinical medicine are methods of imaging anatomical shape changes caused due to diseases. However, upon occurrence of a disease, biochemical changes and molecular genetic changes occur earlier than the biochemical changes.
Therefore, in recent time, a diagnosis paradigm is changing, for example, into monitoring of a quantity of motion, metabolism, biometric information and the like, for diagnosis of a disease in an imaging manner at the beginning of the disease.
Also, wearable devices are currently developed and studied in a direction of providing such devices with a monitoring function for the quantity of motion, the metabolism, the biometric information and the like. In this instance, if one information (object to be measured, property to be measured) is monitored for each sensor, there is a difficulty that a size of the wearable device increases in order to monitor various types of information. Therefore, a method of simultaneously monitoring various types of information through a sensor may be considered.

Measuring apparatus and wearable device comprising same


Embodiments of the present disclosure relate to the communications field, and provide a method for predetermining a resource in random access, user equipment, and a base station, so that a process of random access can be simplified, and reliability of information transmission in the process of random access can be improved. The method is: an implicit manner is used to indicate some or all of identifiers of a random access preamble index; so that a message in a procedure of random access is simplified, and at the same time resources used for transmitting an RAR reply message, transmitting a contention resolution message, retransmitting an RAR reply message, and retransmitting a contention resolution message are configured by using signaling or preset by a system. The embodiments of the present disclosure are used for a process of random access.

In the prior art, an important application of machine-type communication (MTC) user equipment (UE) is a smart meter. A smart meter is usually installed in a basement of a house, or is isolated by a metal casing. The MTC UE is subject to severer a path or penetration loss. Therefore, an operator expects a coverage area of a cell to be further enhanced, so that a service can be provided to MTC UE in an extreme environment (for example, located in a basement or located at an edge of a cell).
During enhancement of a coverage area, signal quality is usually improved by consuming more resources (for example, time resources, and/or frequency resources, and/or power resources, and/or code resources), and/or interference intensity is reduced through interference coordination.
In a wireless communications system, a random access process is usually required to establish a connection between user equipment and a network. An Long Term Evolution (LTE) network is used as an example, and a random access process of LTE has two modes: a contention-based mode and a non-contention-based mode. Contention-based random access includes four steps: Step 1 UE randomly selects one preamble from available preambles (preamble), and sends the randomly selected preamble on one available physical random access channel (PRACH) resource. Step 2 A base station replies to the detected preamble, where the reply message is included in a media access control packet data unit (MAC PDU). The MAC PDU includes one MAC header (header) and one or more random access response (RARs) and a possible padding (padding) bit.
The MAC PDU that carries the RAR may be referred to as a message 2 in short. The message 2 is carried by one physical donwlink shared channel (PDSCH). A physical donwlink control channel (PDCCH) or an enhanced physical donwlink control channel (EPDCCH) schedules transmission of the PDSCH. Therefore, the UE first needs to detect the PDCCH or the EPDCCH, to obtain scheduling information of the PDSCH, and further detect the PDSCH according to the scheduling information. If the UE successfully detects the PDSCH, and a preamble identified by an RAPID included in the MAC PDU carried by the PDSCH is a preamble sent by the UE, the UE believes that the base station has detected a preamble sent by the UE, and believes that an RAR corresponding to a subheader including the RAPID is an RAR sent by the base station to the UE. Step 3 After determining that the RAR fed back by the base station is received, the UE sends an RAR reply on a PUSCH (Physical Uplink Shared Channel, physical uplink shared channel), where the RAR reply is referred to as a h message 3 in short. The message 3 mainly includes a ue-Identity (user identity) and an establishment cause (a cause of an establishment request). Step 4 The base station detects the message 3 sent by the UE, determines that random access of the UE succeeds, and sends a random access contention resolution message to the UE, where the random access contention resolution message is referred to as a message 4 in short. The message 4 is also carried by the PDSCH. The PDCCH or the EPDCCH schedules transmission of the PDSCH. Therefore, the UE first needs to detect the PDCCH or the EPDCCH, to obtain scheduling information of the PDSCH carrying the message 4, and further detect the PDSCH according to the scheduling information. If the UE successfully detects the PDSCH, and the ue-Identity included in the message 4 carried in the detected PDSCH is the same as or matches the ue-Identity, of the UE, included in the message 3, the UE believes that random access succeeds.
It can therefore be seen that in a contention-based random access process, the foregoing four steps involves transmission of sequence, data, and scheduling (or control) information on six channels. A preamble sequence is transmitted on a PRACH channel, the scheduling information of the message 2 is transmitted on the PDCCH or the EPDCCH, the message 2 is transmitted on the PDSCH, the message 3 is transmitted on the PUSCH, the scheduling information of the message 4 is transmitted on the PDCCH or the EPDCCH, and the message 4 is transmitted on the PDSCH. In the foregoing cases, a coverage enhancement requirement is considered. If it needs to be ensured that random access of UE is performed reliably, performance needs to be enhanced for transmission of the foregoing six channels, which means that more resources need to be consumed, and complexity of implementing a base station and UE is greatly increased. Therefore, a mechanism is needed to simplify a random access process, so that implementation of coverage enhancement does not need to consume more resources, thereby reducing complexity of implementing a base station and UE.

Method for predetermining resource in random access, user equipment, and base station


Disclosed is a zoom lens comprising at least three movable lens units. The zoom lens comprises a first lens unit disposed closest to an object side and having positive refractive force. The first lens unit includes a positive cemented lens constituted by a positive lens element and a negative lens element, with a concave side of the cemented surface facing the object side. The zoom lens also comprises a last lens unit disposed closest to an image side. Upon zooming from a wide angle end to a telephoto end, at least the first lens unit and the last lens unit move toward the object side so as to increase an air gap between the first lens unit and a lens unit adjacent to the first lens unit on the image side but to decrease an air gap between the last lens unit and a lens unit adjacent to the last lens unit on the object side. The zoom lens satisfies the following condition: EQU 0.7<f1/(fw.multidot.ft).sup.1/2 <1.4. where f1 is the focal length of the first lens unit, fw is the focal length of the entire lens system at the wide angle end, and ft is the focal length of the entire lens system at the telephoto end.

1. Field of the Invention
The present invention relates generally to a zoom lens and, more particularly, to a zoom lens suitable for use with a lens shutter type camera or the like.
2. Related Background Art
In recent years, a zoom lens has been dominated in the sector of photographing lenses of a lens shutter type camera.
Also, with a higher performance of the zoom lens, a variety of zoom types have been proposed. Further, because of an advancement of technology for a lens barrel in recent years, there have been proposed a variety of zoom types in which high zoom ratio is attained by use of a zoom lens constructed of three or more movable lens units, a so-called multi-unit zoom lens.
Given first is a generality about the multi-unit zoom lens including the three or more movable lens units.
In the multi-unit zoom lens, there increases freedom for selecting a zooming locus of each lens unit upon zooming from a wide angle end to a telephoto end, and, therefore, freedom in terms of correcting an aberration also increases. Further, because lens units changing magnification for zooming are increased, it is possible to attain a higher performance while increasing a zooming ratio. There has hitherto existed a problem in which a structure of a lens barrel becomes intricate with an increment in the number of movable elements. This problem is, however, overcome to some extent owing to the advancement of the technology for the lens barrel in recent years.
Hitherto, in the zoom lens with no restriction of a back focus, a negative lens unit has been disposed closest to an image side in order to reduce the total length of the lens system and to make the lens size small. The zooming is effectively performed by augmenting a variation in the back-focal distance upon zooming from the wide angle end to the telephoto end.
Further, a positive lens unit is disposed closest to the object side in the lens system, thereby attaining the reduction in the total length at the telephoto end.
From the above, a positive/positive/negative 3-unit zoom lens and a positive/negative/positive/negative 4-unit zoom lens may be specifically cited as a multi-unit zoom lens capable of increasing the zoom ratio and suitable for the down-sizing, and a variety of proposals pertaining thereto have been made.
The positive/positive/negative 3-unit zoom lens is constructed of, sequentially from the object side, a first lens unit having refractive power, a second lens unit having the positive refractive power and a third lens unit having negative refractive power. This zoom lens is constructed so as to, upon zooming from the wide angle end to the telephoto end, increase an air gap between the first lens unit and the second lens unit but to decrease an air gap between the second lens unit and the third lens unit (e.g., Japanese Patent Application Laid-Open No. 2-73211).
On the other hand, the positive/negative/positive/negative zoom lens is constructed of, sequentially from the object side, a first lens unit having the positive refractive power, a second lens unit having the negative refractive power, a third lens unit having the positive refractive power and a fourth lens unit having the negative refractive power. This zoom lens is constructed so as to, upon zooming from the wide angle end to the telephoto end, increase an air gap between the first lens unit and the second lens unit but to decrease both an air gap between the second lens unit and the third lens unit and an air gap between the third lens unit and the fourth lens unit (e.g., Japanese Patent Application Laid-Open No. 3-39920).
Also, in the zoom lens with no restriction in terms of the back-focal distance that is used for a lens shutter type camera, it is required for attaining the down-sizing of the camera body to reduce both a size of each lens and a thickness of the lens system (collapse lens thickness) when collapsed in the camera body. In this case, there are reduced both a length (total lens thickness) of the lens system along an optical axis extending from the closest-to-object surface to the closest-to-image surface and a thickness of each lens unit. These reductions are effective in the decrease in the collapse lens thickness and, in turn, the down-sizing of the camera body.
However, in the conventional multi-unit zoom lens such as the positive/positive/negative 3-unit zoom lens and the positive/negative/positive/negative 4-unit zoom lens, the negative lens element of the first lens unit is disposed closest to the object side in order to well correct a positive distortion at the wide angle end and obtain a sufficient back-focal distance. Then, if the first lens unit is constructed of two lens elements, a negative spherical aberration is insufficiently corrected. A well-corrected negative spherical aberration entails nothing but to increase the number of the constructive lens elements or to make the lens surface aspherical. As a result, there arises an inconvenience of being insufficient in terms of simplifying the configuration and reducing the costs.


A power routing device includes processing circuitry configured to receive inputs from device monitors configured to determine load status information from a plurality of loads, power-source monitors configured to determine source status information from a plurality of power sources, and command modules configured to provide context information regarding the loads or the power sources. Some of the loads have different operational characteristics than others. Some of the power sources have different operational characteristics than others. The power routing device is operably coupled to power-source distribution hardware selectively operably coupling selected ones of the power sources to selected ones of the loads. The power routing device is configured to direct the power-source distribution hardware to define routing between the selected ones of the power sources and the selected ones of the loads based on the source status information, the load status information and the context information.

Some mechanical devices, such as actuated prosthetic limbs, include multiple components that require different levels of power in order to operate. In this regard, for example, various ones of the components inside the prosthetic limb may perform different operations and have different ranges of motion. Thus, individual components inside the device may also require varying amounts of power that correspond to different sets of operations. As an example, turning a robotic wrist requires different amounts of power depending upon whether the wrist is merely turning for repositioning purposes, or the wrist is turning to attempt to open a jar. Thus, the mechanical load may differ for different operations requiring essentially the same movement.
In healthy human physiology, the spatiotemporal distribution of power delivery (glucose and oxygen delivered by blood) for limb movements is dynamically controlled based on muscle demands by regulating arterial resistance. However, existing prototypes of actuated prosthetic limbs typically match power demands with a static type or location of source. In other words, a few power sources, each having a fixed power/energy density value, and targeted to power a specific part within the limb, are typically employed without any ability to define dynamic control, resulting in limited power routing capabilities to devices and systems that have complex and variable loading conditions.


Provided is a heat conversion device including a thermoelectric element, the heat conversion device capable of implementing a heat conversion function of high efficiency in spite of low power consumption by disposing a heat absorption module and a heat emission module in a horizontal direction in a structure of a heat exchange device to which a thermoelectric module is applied, and by implementing a heat conversion effect for air while maintaining a desired air volume and air velocity without flow path resistance.

Field of the Invention
Embodiments of the present invention relate to a heat conversion device including a thermoelectric element.
Description of the Related Arts
In general, a thermoelectric element including a thermoelectric device has a structure in which a PN junction pair is formed by bonding a p-type thermoelectric material and an n-type thermoelectric material between metal electrodes. When a temperature difference is provided between the materials of this PN junction pair, electric power is generated by the Seebeck effect so that a thermoelectric element can serve as a power generating device. Also, the thermoelectric element may be used as a temperature controlling device by the Peltier effect in which one material of the PN junction pair is cooled and the other material is heated.
In such a case, the thermoelectric element may be applied as a device that performs both dehumidification and drying of entering air by forming a radiant heat member on a heat emission part and a heat absorption part. However, in this case, even though the air entered from the thermoelectric into the heat absorption part is primarily dehumidified, and the air is then introduced to the heat emission part arranged at an upper portion of the thermoelectric element so as to be dried, it is problematic in that dehumidification efficiency is reduced because the flow of air is increased and flow path resistance is generated. In order to settle this problem, there is a way to entirely make the flow of air strong by increasing the strength of entering air, but this way leads to an increase of power consumption and an increase in noise. Thus, it is problematic in that performance as a dehumidifier is very reduced.


The present invention relates to a method of manufacturing a semiconductor optical device. The present invention discloses a method of manufacturing an optical device of a planar buried heterostructure (PBH) type by which an active layer region of a taper shape at both ends is patterned, an undoped InP layer is selectively grown in order to reduce the propagation loss and two waveguides are simultaneously formed by means of a self-aligned method, thus simplifying the process to increase repeatability and yield.

1. Field of the Invention
The invention relates generally to a method of manufacturing a semiconductor device, and more particularly to, a method of manufacturing a semiconductor optical device of a planar buried heterostructure (hereinafter called xe2x80x98PBHxe2x80x99) type in which a mode converter is integrated.
2. Description of the Prior Art
A semiconductor optical amplifier is usually used as a functional device of a wavelength division multiplexing (WDM) optical communication network. There has recently been made a lot of research and development on the semiconductor optical amplifier. The type of the semiconductor optical amplifier can be functionally classified into a mode converter for reducing the optical loss caused when a semiconductor optical device and an optical fiber are coupled, and one in which current is injected into an active layer for the purpose of the gain of a semiconductor optical device.
The mode converter (or mode expander) is to reduce the optical loss that is caused when the semiconductor optical device and the optical fiber are coupled. The optical loss is generated when a mode of the optical fiber of about 8xcx9c10 xcexcm2 in the mode size and a mode of the optical device of about 1xc3x970.5 xcexcm2 in the mode size due to a large refractive index of the semiconductor is mismatched. The above optical loss requires an optical alignment of a higher degree upon an optical packaging, and also greatly affects the characteristics of an active optical device due to a reflecting light caused in the cross section of the optical device. Therefore, research has recently been made on a mode converter integrated optical device in order to reduce the cost and accomplish single integration between the optical devices.
Next, in case of the one in which current is injected into the active layer for the purpose of the gain of the semiconductor optical device, the two waveguides are different in the width and have a vertical type in space, in the optical device having two waveguides as a basic structure. Therefore, there is a problem in technically regrowing an electrical isolation layer and in repeatability, etc. accordingly.
In order to solve these problems, U.S. Pat. No. 6,025,207 discloses a method of simplifying the process by improving the electrical isolation layer. FIG. 1 is a cross-sectional view of an optical device of a buried ridge structure (hereinafter, called xe2x80x98BRSxe2x80x99) disclosed in the U.S. Pat. No. 6,025,207.
Referring now to FIG. 1, the optical device of BRS type in the U.S. Pat. No. 6,025,207 includes, as basic components, an optical waveguide layer 102 for integration and an active layer 104. A p-InP layer 106 for current injection is regrown on the entire structure. A proton 108 is injected into regions except for the active layer 104 using the ion implanter and current is injected into the active layer 104 in order to simplify the process for integration. Further, an InGaAs layer 110 into which a p type impurity is doped and having a high conductivity as an ohmic layer is formed. An electrode 112 is formed to complete the optical device. However, as there exists the difference in the width of current injection and in the width of the active layer, electrons are leaked to the p-InP layer 106 when high current is injected. As the InP layer 106 into which a p type impurity is doped is regrown on the first waveguide in regions except for the active layer 104, the propagation loss occurs. There is a problem that the propagation loss degrades the characteristic of an active optical device. In FIG. 1, unexplained reference numerals 100, 101 and 103 indicate a substrate, a buffer layer, and a medium layer into which a n type impurity of a high concentration is doped, respectively.
U.S. Pat. No. 5,863,809 discloses a method of reducing this propagation loss. FIG. 2 is a cross-sectional view of an optical device of a buried ridge structure (BRS) type disclosed in the U.S. Pat. No. 5,863,809.
Referring now to FIG. 2, waveguide layers 202 and 204, and an active layer 206 are first grown. A waveguide of the active layer 206 is then formed using photolithography and etching. An undoped layer 210 is regrown on the entire surface to reduce the propagation loss.
Further, the undoped layer 210 on the active layer 206 is etched by means of photolithography so that current can be injected into only the active layer 206. InP layers 214 and 212 being an electrical isolation layer and being doped by p/n/p type, are formed in regions except for the active layer 206. The U.S. Pat. No. 5,863,809 can obtained a current isolation characteristic using the undoped InP layer 210. However, the patent has a disadvantage that current is distributed due to the thickness of the undoped layer 210. In addition, as photolithography and etching process are performed several times, repeatability is degraded and the yield is thus reduced. Considering that the yield of the optical devices is very lower than that of electronic devices, reduction in the yield must be solved. In FIG. 2, unexplained reference numerals 207, 208, 216, 218 and 220 indicate a p-InP layer, an etch stop layer, an upper p-InP layer, a contact layer and a contact pad, respectively.
The present invention is contrived to solve the above problems and an object of the present invention is to provide a method of manufacturing a semiconductor optical device having a PBH (planar buried heterostructure) type of a good characteristic, capable of reducing the propagation loss, simplifying the process using a self-aligned method, and increasing the repeatability and yield.
In order to accomplish the above object, a method of manufacturing a semiconductor optical device according to the present invention, is characterized in that it comprises the steps of sequentially forming a first waveguide layer, a first clad layer, a second waveguide layer and a second clad layer on a semiconductor substrate; depositing a first hard mask on the second clad layer and then forming a first hard mask pattern having a taper shape at both ends; etching the second clad layer and the second waveguide layer using the first hard mask pattern as an etch mask; selectively growing an undoped InP layer at regions from which the second clad layer and the second waveguide layer are etched using the first hard mask pattern as a selective mask and then patterning the undoped InP layer; removing the first hard mask pattern; forming a second hard mask pattern the width of which is smaller than the width of the first hard mask pattern on a result from which the first hard mask is removed; and etching the undoped InP layer, the second clad layer, the second waveguide layer, the first clad layer and the first waveguide layer using the second hard mask pattern as an etch mask to simultaneously form a second waveguide layer and a first waveguide layer having different widths.
The method of manufacturing a semiconductor optical device according to the present invention, is characterized in that it can further comprises the steps of, after the step of etching the undoped InP layer, the second clad layer, the second waveguide layer, the first clad layer and the first waveguide layer using the second hard mask pattern as an etch mask, forming an electrical isolation layer at regions where the undoped InP layer, the second clad layer, the second waveguide layer, the first clad layer and the first waveguide layer are etched; removing the second hard mask pattern; forming an upper clad layer on a result from which the second hard mask pattern is removed; depositing an electrically conductive layer on the upper clad layer and then patterning the electrically conductive layer in order to inject current into the second waveguide layer; forming a silicon nitride film on a result in which the electrically conductive layer is pattern and then etching the silicon nitride film in order to electrically connect the electrically conductive layer and an electrode, thus forming an opening portion; and forming an electrode electrically connected to the electrically conductive layer.
The electrical isolation layer is formed by growing an InP layer in a p/n/p type.
The first waveguide layer is formed to be a n-InGaAsP layer, the second waveguide layer is formed to be an undoped InGaAsP layer, the first clad layer is formed to be a n-InP layer, and the second clad layer is formed to be a p-InP layer.
An InGaAsP layer is further formed at lower and upper sides of the second waveguide layer for separate confinement heterostructure and uniform gain depending on polarization.
The widths at the central portion and edge of the second hard mask pattern are different, the width at the edge of the second mask pattern is wider than the width at the central portion of the second mask pattern, the width of the central portion of the second mask pattern is narrower than the wide of the fist hard mask and the boundary at the central portion and the edge of the second mask pattern is tapered
The etching process using the second hard mask pattern as an etch mask includes first performing a dry etching process to etch by a given depth and then performing a wet etching process.
The first and second hard masks are formed to be a silicon nitride film.

Method of manufacturing semiconductor optical device

Electronics and Telecommunications Research Institute


A method of coating a surface, preparing a doped metal-ion precursor solution for coating, and an article including a component coated by the described method are disclosed. The method of coating includes applying a fluoro-silane doped metal-ion precursor solution on the surface to form a coated surface. The metal-ion precursor solution includes greater than about 0.6 molar percent concentration of a metal-ion precursor in a solvent comprising an alcohol. The method of preparing the doped metal-ion precursor solution includes dissolving a metal-ion precursor in a solvent comprising an alcohol at a temperature greater than about 100° C. and refluxing at a temperature greater than about 150° C. such that the concentration of metal-ion precursor in the solution is greater than 0.6 molar percent of the solution, and adding a fluoro-silane to the metal-ion precursor solution.

The invention relates generally to a method of coating a surface and an article incorporating the coated surface. More particularly the invention relates to coating a surface with a doped metal-ion precursor solution and making a coated surface having a doped metal-ion compound.
Hydrophobic and super-hydrophobic surfaces are desirable in numerous applications, such as windows, DVD disks, cooking utensils, clothing, medical instruments, automotive and aircraft parts, textiles, and like applications. Typically hydrophobic surfaces have been created by changing surface chemistry or by surface texturing, or by combining both of these methods. Altering the surface chemistry of the surface in order to achieve hydrophobicity typically involves coating the surface with a low surface energy coating. However, most of such hydrophobic coatings suffer from lack of durability in their application environment with issues of poor adhesion to the surface, low mechanical robustness and lost functionality. Moreover, most of the existing techniques for altering the wetting resistance of the surface suffer from certain drawbacks, such as processes that are time consuming, difficult to control, expensive or ineffective in producing films with sufficient durability, and are usually not scalable to large surface areas. Therefore, there is a need for an inexpensive, easy, and effective means for achieving surfaces with wetting resistance.

Method of coating a surface and article incorporating coated surface


An arbiter connects a plurality of devices to a bus. The arbiter determines priority among the devices based on minimum access intervals associated with the devices and timers which keep track of the elapsed time since each device last had access to the bus. The timers can be configured to either count up from zero to the minimum access interval of each device or count down from the minimum access interval of each device to zero. The arbiter can also adjust the minimum access intervals of the various devices based upon factors such as the amount of data required by the device, the amount of data most recently received by the device and the transfer rate of the device. The arbiter thus optimizes bus usage while minimizing the likelihood of a given device not functioning based on an ability to access the bus due to contention with other devices.

When many devices are connected to a data bus, such as in a multi-media application, conflicts will occur where two devices request access to the bus at the same time. To resolve these conflicts, arbiters have been developed to grant bus access to the device having the highest priority. Thus, when two devices are in contention for the bus, the arbiter first gives access to the device with the higher priority of the two and subsequently gives access to the device with the lower priority.
In a system with such an arbiter, a device with a very high priority will get frequent access to the bus while a device with a low priority will get infrequent access to the bus. This can present a problem when the devices are assigned their priorities during system initialization or boot-up such that the relative priorities of the devices do not change over time. A situation can exist where a low priority device needs access to the bus in order to continue functioning but cannot obtain access because a high priority device which does not necessarily need access happens to be requesting bus access coincidentally. Thus, although the system is granting access on a priority basis, inefficiencies can result where low priority devices experience idle time or malfunction.

Multimedia system employing timers to properly allocate bus access


A recognition system for recognizing a data packet includes a data memory for storing data units in a predetermined order. A first data unit includes an instruction for determining a characteristic of the data packet. The first data unit also includes the total number of the data units, and a corresponding sequence number for the first data unit. The recognition system further includes a matching processor coupled to the data memory. The matching processor can execute the instruction in the first data unit to generate a result. The matching processor can also provide an address pointer according to the execution result, the total number of the data units, and the corresponding sequence number of the first data unit.

Conventional protocol recognition systems can determine an application layer protocol (e.g., hypertext transfer protocol or HTTP, file transfer protocol or FTP, e-donkey protocol) of IP (Internet Protocol) packets based on a pattern matching process. For example, a conventional protocol recognition system can include a universal computer with pattern matching software. The universal computer can use the pattern matching software to compare the content of the IP packets with a set of patterns in the pattern matching software, so as to determine the protocol of the IP packets according to the result of the comparison. However, such conventional protocol recognition systems may not be able to determine the protocol of the IP packets fast enough.
In another example, a conventional protocol recognition system can be realized by a hardware circuit in an FPGA (field-programmable gate array). For example, an HDL (hardware description language) program such as a Verilog program for determining the protocol of the IP packets can be compiled to the hardware circuit. As such, the hardware circuit can process the IP packets so as to determine the protocol of the IP packets according to the HDL program. However, usually a relatively large FPGA is configured to load the hardware circuit, which can increase the size and the cost of the recognition system. Additionally, it would take a relatively long time to write the HDL program due to the complicated design of the hardware circuit. Moreover, any change to the patterns for the application layer protocol (e.g., HTTP, FTP, e-donkey protocol, BITTORRENT protocol, simple mail transfer protocol or SMTP), or any new protocol added to the Internet protocol suite, will need a programmer to re-write the HDL program, which would also take a relatively long time.

Recognition systems based on pattern matching