A monitor device and associated methodology are disclosed which provide a self contained, relatively small and continuously wearable package for the monitoring of heart related parameters, including ECG. The detection of heart related parameters is predicated on the location of inequipotential signals located within regions of the human body conventionally defined as equivalent for the purpose of detection of heart related electrical activity, such as on single limbs. Amplification, filtering and processing methods and apparatus are described in conjunction with analytical tools for beat detection and display.

The heart is a muscular pump that is controlled by a natural electrical system that causes the heart muscle to contract and pump blood through the heart to the lungs and the rest of the body, carrying oxygen as well as other needed nutrients. The heart can be characterized by a set of parameters that describe the state of the heart, including the frequency and timing of the contractions of the four chambers of the heart, and the pattern of electrical signals causing those contractions. There are many methods of detecting these parameters that are well known in the art, including: sensing the electrical impulses of the heart, sensing the pulse of blood as it moves through arteries, Doppler and other acoustic based methods, capacitance, micro-impulse radar, pressure- and/or motion-based methods such as by utilizing piezo-electric elements or strain gauges, and optical methods in areas where the pulsing of blood can be externally viewed, such as in a pulse-oximeter.
The most well-known and conventional method utilized today for measuring heart-related parameters is the electrocardiogram. An electrocardiogram, or ECG, signal is a surface measurement of the electrical potential of the heart generated by electrical activity in cardiac tissue. This measurement can be made using electrodes placed on the surface of the skin because the entire body is capable of conducting electricity.
FIG. 1 shows a typical ECG signal generated by one heart beat. Signal strength is shown on the Y axis and time is shown on the X axis. The individual spikes and dips in the signal are called waves. The P wave shown in FIG. 1 represents the contraction of the atria. The Q, R, and S waves, referred to as the QRS complex, represent the contraction of the ventricles. The T wave represents the recovery, or repolarization, of the ventricles. The amplitude of a typical ECG signal is approximately 1 to 2 mV when measured from the chest using good electrode contacts.
ECG measurements may be used to provide information about a number of heart related parameters, including, but not limited to, the heart beat rate, or heart rate, for a number of applications, such as medical diagnostic, health awareness and sports performance applications. The most reliable heart rate calculation based upon ECG is performed by detecting each QRS complex, and thus each heart beat, because the QRS complex contains the highest amount of energy and its spectrum differs sufficiently from the spectrum of movement artifacts. Beats are typically counted at each R point (the peak), and the distance between a first R point and a subsequent R point is known as the R-R interval, which, when inverted, yields the instantaneous heart rate. Other parameters such as the heart-rate variability are also computable from the set of R-R intervals.
As discussed above, the heart is a source of a voltage potential difference resulting from the electrical activity that causes the heart muscles to contract. This potential difference is known in the art as the heart's action potential. An ECG signal is a measurement of this action potential. In addition, the heart is positioned in the left chest area and is oriented at an angle slightly off of vertical. The traditional model of ECG measurement indicates that ECG measurements must be taken across the heart, meaning using electrodes placed on either side of an imaginary line running through the center of the heart. Many different researchers have identified the various sections of the surface of the body in different ways with respect to placing electrodes for measuring different aspects of the heart's electrical activity.
Generally, these placements are identified in two ways. First, pairs of electrodes are often used to measure the electrical potential difference between two points. If two points show an electrical potential signal that varies with the activity of the heart they are said to be not equipotential or therefore inequipotential with respect to one another. Inequipotential therefore refers solely to the difference in the heart's action potential rather than other sources of voltage difference such as EMG. Furthermore, locations are described herein as measuring a different aspect of the heart's electrical activity from other locations when those two locations are inequipotential. The electrodes are conventionally placed in a manner as to obtain maximum differentiation between the electrodes. Conventionally, therefore, the body is divided into quadrants I, II, III and IV, as illustrated in FIG. 1A. Electrodes are conventionally placed in two different quadrants on the body, where the body 1 is divided into four sections, or quadrants, by two planes running through the heart. The location of these planes has been modified over time as knowledge in this field has progressed, but has remained fairly constant in that sagittal plane 2 runs roughly vertically through the heart and the transverse plane 3 roughly horizontally. These two planes are orthogonal to one another when viewed from the two-dimensional perspective from in front of the patient. It is important, for the purposes of this application, to assess the location of these imaginary planes through the heart. Sagittal plane 2, is sometimes considered to be coincident with the medial line of the body. Other views, however, direct the vector along a more canted axis coincident with the slightly asymmetric orientation of the heart within the chest cavity. Transverse plane 3, is orthogonal to sagittal plane 2. For bipolar electrode placements, the two electrodes are conventionally placed in two different quadrants, allowing the measurement of the heart's action potential. The other method of reading ECG signals from the heart is to take single pole readings that utilize a single electrode at one point and then utilize an average of multiple electrodes for the other point. This allows a view of the heart from different directions, and allows the creation of views of the heart not achievable with only two electrodes. The precordial, or chest placements in the standard 12-lead ECG are examples of this sort of placement.
Other models include the Einthoven triangle, which describes a roughly inverted equilateral triangular region on the chest having a base extending between the left and right shoulder joints and an apex approximately located at the base of the ribcage, below the sternum. The model contemplates the angle formed at the right shoulder having a first aspect of the ECG signal, the abdominal angle having a second such aspect and the left shoulder angle having a third aspect. The Bayley triaxial system and the Hexaxial system each divide the chest and abdominal area into a larger number of sections or regions, each of which is assigned a single aspect or mixed aspect of the ECG signal.
All of the prior art location identification systems require electrodes placed in at least two of the quadrants of the body. The surface area of each quadrant is defined herein, therefore, as an equivalence region on the body, the portions of the body near the boundaries of the quadrants are further eliminated from such equivalence regions, as it is commonly understood that the boundary can move slightly as the heart beats, the person moves, and that the boundaries can be different between different individuals due to minor difference in heart orientation within the body. The equivalence regions are thus defined as the quadrants illustrated in FIG. 1A. Previous systems for measuring ECG all require having electrodes in at least two of the equivalence regions. These equivalence regions as well as the plethora of different mappings applied to the surface of the body utilized in the prior art can also be understood as following the principle that the signals obtainable within these quadrants are homogeneous as it is assumed that the body is composed of a homogeneous material.
Several prior art devices exist for measuring ECG based on the traditional model. For example, clinical or medical ECG devices use several electrodes placed on the chest, arms and legs to measure a number of different ECG signals from selected electrode pairs wherein in each pair, one electrode is located in one equivalence region and the other electrode is located in a different equivalence region. The different readings together allow a clinician to get a view of the function of the three dimensional electrical activity of the heart from a number of different angles. In many cases, the devices which provide the ability to detect and monitor the heart related parameters is stationary and is intended to monitor a stationary patient.
Such devices, while highly accurate, are very expensive and cumbersome and thus do not lend themselves well to ambulatory or long term uses such as in a free living environment. Holter monitors are devices that may be used for continuous, ambulatory ECG measurement, typically over a 24-48 hour time period. These Holter devices collect raw electrical data according to a preset schedule, or frequency, typically 128 hz or 256 hz. These devices must therefore contain a significant amount of memory and/or recording media in order to collect this data. The physical bulk and inconvenient accessories of this device restricts its continuous use to a relatively short time frame. Each device comprises at least two electrodes for clinical or monitoring data detection and typically a third electrode for ground. The leads are designed to be attached to the chest across the heart, or at least across the conventionally understood sagittal plane 2, and a monitoring device connected to the electrodes is carried or worn by the patient, which is typically a heavy rectangular box clipped to the patient's waist or placed in a shoulder bag. The sensors utilized in conjunction with the device are affixed according to a clinical procedure, wherein the skin under the electrode or sensor is shaved and/or sanded and cleaned with skin preparation liquids such as alcohol prior to application to improve signal quality. Consequently, the sensors are not easily interchanged and may limit physical or hygienic activity. Holter monitors are relatively expensive and for the reasons listed above, are not comfortable for long term and/or active wear situations.
Loop monitors are configured and worn similarly, yet are designed to work for longer periods of time. These systems are designed to record shorter segments or loops of raw data or morphology when the wearer signifies, by pressing a time stamp button that they are doing an activity of interest or feeling a chest or heart related pain. The device will typically record 30seconds before and 30 seconds after the time stamp. While some success with respect to longer term wearability and comfort is achieved, these loop monitor devices are still inconvenient for everyday use, and include lead wires from the device, snap on sensors affixed to the body by adhesives which require daily skin preparation and periodic re-alignment of the sensors to the original positions.
More recently, a few monitors have also been provided with some automated features to allow the device, without human intervention, to record certain loops when certain preset conditions or measurement thresholds are achieved by the detected heart related activity, such as an abnormal beat to beat interval or a spike in heart rate. Implantable loop recorders have also been developed, which provide similar functionality, with the attendant inconvenience and risks associated with an invasive implant.
Another diagnostic device is known as an event recorder, and this device is a hand held product, with two electrodes on the back, some desired distance apart with recording capabilities where a patient is instructed to place this device against the skin, over the heart or across the sides of the body in order to record a segment of data when the patient is feeling a heart related symptom. This device is not utilized for continuous monitoring, and has memory capability for only a limited number of event records. Once the media storage is filled, there is a facility on the device to communicate the data back to a clinic, clinician, service, or doctor for their analysis, usually by telephone.
While not designed for medial or clinical applications per se, a number of chest strap heart rate monitors have been developed that may be used to measure heart rate from ECG, with some recent devices being capable of recording each detected heart beat, recorded in conjunction with a time stamp in the data. Examples of such conventional monitors commercially available include Polar Electro Oy, located in Oulu, Finland and Acumen, Inc. located in Sterling, Va. These chest strap monitors are designed to be wrapped around the torso beneath the chest and include two electrodes positioned on either side of the heart's conventionally understood transverse plane 3 for measuring an ECG signal. The device is placed just below the pectorals, with conventional electrode positioning. The device is placed at this location because noise and motion signal artifacts from muscle activity is minimal and the signal amplitude is quite robust, consistent and discernable by a circuit or software application. Chest strap monitors of this type, while promoted for use in exercise situations, are not particularly comfortable to wear and are prone to lift off of the body during use, particularly when the wearer lies on his or her back.
Finally, a number of watch-type ECG based heart rate monitors are commercially available, such as the MIO watch sold by Physi-Cal Enterprises LP, located in Vancouver, British Columbia. Such watches include a first electrode attached to the back of the watch that, when worn, contacts one arm of the wearer, and one or more second electrodes provided on the front surface of the watch. To get an ECG signal, and thus a heart rate, a wearer must touch the second electrode(s) with a finger or fingers on the opposite hand, that is, the hand of the arm not wearing the watch. Thus, despite being worn on one arm, the watch measures ECG according to the conventional method, being across the heart, again on either side of the heart's conventionally understood sagittal plane 2, because the two electrodes are contacting both arms. Such watches, while comfortable to wear, only make measurements when touched in this particular manner and thus are not suitable for monitoring ECG and heart rate continuously over long periods of time or while conducting everyday activities such as eating, sleeping, exercising or even keyboarding at a computer.
Matsumara, U.S. Pat. No. 5,050,612, issued Sep. 24, 1991, discloses the use of a multi-electrode sensing watch device, identified as the HeartWatch, manufactured by Computer Instruments Corporation, Hampstead, N.Y., for certain types of heart parameter detection. While Matsumara discloses that the conventional use of the HeartWatch device is in conjunction with a chest strap, he also identifies an alternative use which relies solely on the multisensor watch device itself. The device has two electrodes at different distances along the arm from the heart, and the detected waveform from one electrode is subtracted from the other to obtain a resultant signal. Matsumura identifies this signal as not resembling an ECG, but states that it is useful for detecting ST segment depression. No teaching or suggestion of the efficacy of this method for the identification of heart rate or other heart related parameters is made.
As described above, the traditional models of ECG measurement do not contemplate the action potential of the heart, and thus ECG, being detected and measured from two points within a single quadrant or within a single equivalence region. Moreover, the traditional model rejects the measurement of the action potential from two locations on the same limb. The prior art does contemplate some sensor placements which take advantage of the three dimensional nature of the human body and allow for measuring the heart's action potential between electrodes placed on the front and back of the body, or between spots high on the torso and low on the torso, but on the same side of the body. One skilled in the art would recognize that the prior art only utilized sensor placements that included two or more electrodes in multiple quadrants or equivalence regions.
Another significant shortcoming of ambulatory devices is electrical noise. Noise is detected from both ambient sources surrounding the body, movement and organ noise within the body, and most significantly, the movement of the body itself, including muscle artifacts, motion artifacts, skin stretching and motion between the electrode and the skin. A variety of patents and other references relate to the filtering of noise in many systems, including heart rate detection. In Zahorian, et al., U.S. Pat. No. 5,524,631, issued Jun. 11, 1996, a system is disclosed for detecting fetal heart rates. A significant noise problem exists in that environment, including the heart action of the mother, as well as the significant noise and distortion caused by the fetus'location within a liquid sac inside the mother's abdomen. Zahorian utilizes multiple parallel non linear filtering to eliminate such noise and distortion in order to reveal the fetus' heart rate. The system, like many of the prior art, is unconcerned with the wearability of the monitoring device or the ability to continuously monitor the subject over a long period of time.
None of the above systems identified above combine wearability and accuracy in a compact device. What is lacking in the art, therefore, is a device which provides the ability to measure ECG from two locations in a single equivalence region, such as within a single quadrant as shown in FIG. 1A or on a single limb. Although there are some examples in the prior art that recognize the possibility of inequipotential pairs within a single equivalence region, that the teachings of the prior art fail to utilize these pairs for obtaining a viable signal. There are several barriers to the ability to utilize these signals from unconventional locations, including the small amplitude of the signal, which can be less than one tenth of the signal measured at most conventionally placed electrode locations, the high amount of noise with respect to that signal, as well as the significant effort and risk required to overcome limitations in accuracy, amplitude, and noise obtained from unconventional placements. What is further lacking in the art is such a device which is relatively small in size and adapted for longer periods of continuous wear and monitoring, in conjunction with sensors which minimize the requirement of clinical observation, application or preparation. Such a device provides new opportunities for continuous heart monitoring, including improved comfort, less complex products, and improved compliance with monitoring. Additionally, what is lacking in the art is the ability to combine the continuous monitoring of the heart related parameters with a device which can detect, identify and record the physical activities of the wearer and correlate the same to the heart related parameters.

Method and apparatus for measuring heart related parameters


There is described a device (1) of the reversible type for the connection and reconstruction of a ligament, comprising a central reconstruction portion (9) and two connecting portions (10a, 10b) at the two ends of said central portion (9), characterized in that said connection device is made of silicone, said central portion covers said two connecting portions (10a, 10b), said two connecting portions (10a, 10b) being generally grid-shaped for receiving fixing means adapted to fix said device (1) to said ligament, said device (1) further comprising an outer carbon coating (14).

Ligaments are formations of fibrous connective tissue and combine two or more anatomical structures. The cruciate ligaments of the knee joint are fibrous bundles crossed with one another to support the joint itself and allow its joint movements. Many individuals—in particular those who practice a sport—undergo damage to these ligaments. These ligaments are not the only ones to get damaged, however they are the most frequent and most critical ones as regards their recovery process.
In order to reconstruct the tendons damaged upon a lesion there exists a technique known as "autologous". According to this surgical technique, a fiber bundle is used, which is taken from the ligaments which are used as a basis to reconstruct the damaged ligament. This technique is very invasive and cannot be always applied.
In response to the "autologous" technique, devices have been developed to "simulate" a ligament. Such devices have the shape and size of the ligament to be replaced and are made of polyester. At the ends thereof, there are fixing means (staples, anchors or absorbable or non-absorbable synthesis screws) to stably fix the artificial ligament and allow the partial reconstruction of the ligament tissue thereon. The implant of artificial ligaments of this type is final, over time a part of ligament reconstructs thereon and integrates its functions. Such devices ensure a good biocompatibility, however they have some drawbacks related to their mechanical function. In fact, their flexibility is acceptable but not excellent. Therefore, the performance of such devices is valid even if with some drawbacks, especially as regards the mechanical functionality. However, such a drawback can be especially insidious because such devices can be damaged even shortly after their installation, thus requiring a new operation to replace the device.
These drawbacks caused the "autologous" technique to be the predominant choice despite its limitations.

Non-absorbable tissue reconstruction device, in particular for tissues such as ligaments


A gas production apparatus including: an electrolysis vessel; first and second electrolyte circulation systems; and an electrolyte exchanger, the first/second electrolyte circulation system including: a first/second circulation tank receiving and storing a first/second electrolyte flowing out from an anode chamber/a cathode chamber; and a first/second circulation pump supplying the first/second electrolyte to the anode chamber/the cathode chamber, the electrolyte exchanger transferring part of the first electrolyte existing in the first electrolyte circulation system into the second electrolyte circulation system on one hand, and transferring part of the second electrolyte existing in the second electrolyte circulation system into the first electrolyte circulation system on the other hand.

The alkaline water electrolysis method is known as a method for producing hydrogen gas and oxygen gas. In the alkaline water electrolysis method, water is electrolyzed using a basic aqueous solution (alkaline water) containing a dissolved alkali metal hydroxide (such as NaOH and KOH) as an electrolyte, to generate hydrogen gas at a cathode and oxygen gas at an anode. An electrolytic cell for alkaline water electrolysis usually includes an anode chamber and a cathode chamber separated by an ion-permeable separating membrane, and water is electrolyzed as the electrolyte is circulated through each of the anode and cathode chambers. The electrolyte recovered from each chamber is once collected and stored in a circulation tank, and the electrolyte stored in the circulation tank is supplied to each chamber again.

Gas production apparatus and gas production method


A slow-release irrigation device having a housing with an open and closed end, and a water reservoir located therein, a circumferential channel located proximal to the open end providing a watertight engagement with a neck opening of a beverage container, a water channel connected to the water reservoir and having a water receiving inlet and an outlet for directing water into the water reservoir while simultaneously operating as a water seal for the housing, a housing water outlet located between the ends of the housing and directly connecting the water reservoir with the external environment, an air channel located within the housing, an air tube connecting the water reservoir to an interior of the container; and a permeable membrane located within the air channel and controlling the passage of air through an air outlet of the air tube.

It has been proven that slow delivery of water to plants to match the rate that they can absorb enables growth with very small quantities, by minimizing the amount that is normally lost to lost to run off and absorption into the soil. Furthermore, continuous availability of water maximizes plant growth by allowing photosynthesis to take place whenever light is available. The present invention attempts to deliver both slow delivery and continuous availability of water to maximize the growth, using a minimum of water, while utilizing used beverage bottles which would otherwise be waste.

Slow-release irrigation device for individual plants


An apparatus for purifying air which may act downstream supplement to a catalytic converter of an internal combustion engine for further reducing the hydrocarbon and carbon monoxide gases which the catalytic converter failed to remove. The apparatus comprises a housing filled with volcanic ash over which is directed a turbulent stream of exhaust gases for chemical reaction with the volcanic ash.

This invention relates to the treatment of hydrocarbon and carbon monoxide gases issuing from various sources one of which is the gases issuing from a catalytic converter connected to the exhaust of an internal combustion engine. More particularly, this invention is directed to apparatus positioned downstream of the catalytic converter of an internal combustion engine for supplementing the function of the catalytic converter and may be positioned in the gas stream before, in place of or after the muffler.
Current automotive catalytic converters comprise a multichannel substrate having a flow-through channels which are coated with noble metal catalyst materials that convert the pollutant exhaust gases to harmless effluent.
The automotive catalytic converters generally comprise a large, oval-shaped ceramic substrate housed in a metal sleeve or can. The converters connect directly to a manifold exhaust pipe of an internal combustion engine and normally discharges purified gases through a muffler to atmosphere.
The catalytic converter is designed to reduce hydrocarbon (HC) and carbon monoxide (CO) emissions into carbon dioxide and water vapor. The converter also purges oxides of nitrogen (NOx) vapors from the exhaust. It is called catalytic because precious-metal catalysts inside the converter react with these toxic gases on the large surface area of an ingenious honeycomb coated with less than a gram of the catalytic material.
Most vehicles made after 1980 use this so-called monolithic converter. Pre 1980 vintage cars have a pellet-type converter. Pellet converters contain thousands of tiny grains of ceramic reactants that produce a similar result.
Although the catalytic converter has been successful in removing some of the hydrocarbons and carbon monoxide from the exhaust gases of internal combustion engines, there is still a need for removing still more of these toxic gases exiting from the exhaust of the catalytic converter as well as from other sources of air contamination. This can be accomplished with regard to the catalytic converter by uniformly distributing the exhaust gas flow at a point before, after or in place of the muffler to a new add on supplement converter that agitates the flowing exhaust stream of gases and washes it over a mineral substance to further convert the remaining hydrocarbon and carbon monoxide emissions into carbon dioxide and water vapor.

Apparatus for treating hydrocarbon and carbon monoxide gases


A semiconductor device includes a metal substrate including a metal base plate, an insulating sheet located on the metal base plate, and a wiring pattern located on the insulating sheet, and a semiconductor element located on the metal substrate. The semiconductor element is sealed with a molding resin. The molding resin extends to side surfaces of the metal substrate. On the side surfaces of the metal substrate, the insulating sheet and the wiring pattern are not exposed from the molding resin, whereas the metal base plate includes a projecting portion exposed from the molding resin.

1. Field of the Invention
The present invention relates to the structure of a semiconductor device and a method of manufacturing the same, and more particularly, to a semiconductor device manufactured by a batch transfer molding process.
2. Description of the Background Art
As the method of manufacturing a power semiconductor device, a batch transfer molding process is known in which a plurality of resin-sealed semiconductor devices are collectively formed and then cut (diced) to obtain individual resin-sealed semiconductor devices. The batch transfer molding process enables to improve the reliability of a power semiconductor device and reduce the size and cost thereof.
For example, Japanese Patent Application Laid-Open No. 2002-110885 discloses the method of manufacturing semiconductor devices having a lead frame structure by a batch transfer molding process. Japanese Patent Application Laid-Open No. 2002-110885 describes the technology of providing line-shaped trenches along the dicing lines in the lead frame before cutting a plurality of semiconductor devices that are collectively resin-sealed into pieces through dicing (singulation). This technology lowers the heat generated due to friction in dicing and reduces the occurrence of "burrs" on the cut surface of the lead frame.
In Japanese Patent Application Laid-Open No. 2002-110885, the line-shaped trenches are provided in the lead frame for obtaining advantages in processing, such as ease of cutting and improvement in processing accuracy. According to Japanese Patent Application Laid-Open No. 2002-110885, as shown in FIG. 1, terminal portions cut out from a lead frame are exposed from a molding resin on end surfaces (cut surfaces) of the singulated semiconductor device. Normally, in a semiconductor device having a lead frame structure, a molding resin serves exclusively to insulate terminal portions. Accordingly, with such a configuration, when moisture infiltrates from an interface between the terminal portion exposed at the edge of the semiconductor device and the molding resin, the insulating property of the molding resin may be degraded. The interface peels off in some cases, which may make the package open.
In particular a power semiconductor device is required to withstand high voltage as well as high current, and thus, the reliability of the power semiconductor device may be degraded with the structure in which terminal portions are exposed on the end surfaces of the semiconductor device.

Semiconductor device and method of manufacturing the same


A spherical joint features a housing, a lower race seated therein, a ball member disposed between said lower race and an upper race position, and a retaining ring installed directly over the upper race in an annular groove inside the housing. In one embodiment, a retaining member is threaded into the housing over the retaining ring. The joint is fully serviceable by removal of the retaining member and ring. Prior to component replacement, initial wear can be compensated for, as the groove is oversized so as to originally lie partially below the ring-abutted upper surface of the upper race. When the upper race moves down the body as a result of surface wear of one or more of the races and ball member, tightening of the retaining member once again forces the retaining ring down against the upper race to re-tighten the races against the ball member disposed therebetween.

Conventional tie rod ends are manufactured in a manner that prevents access to interior components for service of the tie rod ends when significantly worn after extended use. In such a conventional structure, the housing of the unit has an open bottom end of reduced diameter presenting an inward projecting lip or flange against which the lower race is seated around the opening through which the stem or stud of the ball member projects. The ball portion of the ball member is received between the lower race and a corresponding upper race, which together provide spherically contoured surfaces conforming against the ball at positions thereover and thereunder so that the ball can rotate in any of three orthogonal axes to allow rotation of the stem and tilting thereof in any direction. In this conventional structure, the top end of the housing is permanently closed off, for example by pressing a cap into the body over a crush-washer surmounting the upper race to crush the washer to a level providing a desired degree of tightness of the upper race against the ball, and then crushing a retaining lip into position over the perimeter of the closed cap. The housing is thus permanently closed at both ends, preventing access to the internal components. Accordingly, the internal components cannot be replaced when worn, and so the entire assembly must be replaced with a new joint, as the internal components cannot be accessed, at least not in a manner that is non-destructive to the housing to allow re-use thereof.
Some serviceable designs for tie rod ends and ball joints have been proposed.
U.S. Pat. No. 6,908,251 teaches tie rod ends and ball joints where an externally threaded retainer member is engaged into the housing from the upper end thereof to close the housing and tighten the upper race against the ball member. The threaded engagement of this retainer member to the housing allows removal for access to the internal components. However, there is some concern that even with use of a set screw to secure the retainer member in place, the threaded coupling may come loose during use of the tie rod end or ball joint.
U.S. Patent Application Publication No. 2009/0238636 teaches tie rod ends and ball joints that likewise use a threaded retainer member, but additionally add a retaining ring and snap ring engaged to the housing overtop of the threaded retainer member. While the design does allow access to internal components for service, and prevents backing off of the threaded retainer member by way of the added rings, there still remains room for improvement.
Applicant of the present application has developed novel spherical joints useful to provide reliable and serviceable rod ends, tie rod ends and ball joints.

Serviceable spherical joints with wear-compensation features


A mounting assembly for mounting a composite pressure vessel to a vehicle includes a saddle having a curved surface extending between two pillars for receiving the vessel. The saddle also has flanged portions which can be bolted to the vehicle. Each of the pillars has hole in which is mounted the shaft portion of an attachment member. A resilient member is disposed between each of the shaft portions and the holes and loaded by a tightening nut. External to the holes, each of the attachment members has a head portion to which a steel band is attached. The steel band circumscribes the vessel and translates the load on the springs into a clamping force on the vessel. As the vessel expands and contracts, the resilient members expand and contract so that the clamping force applied by the band to the vessel remains constant.

In many aircraft and spacecraft the flight control surface, (flaps), are positioned by pneumatic actuators. A pneumatic actuator is an actuator that is operates by high pressure gas. The high pressure gas is commonly stored in steel bottles that are directly bolted to some structure on the aircraft or spacecraft. A disadvantage to using these steel vessels is their weight.
Because keeping the weight of the spacecraft as low as possible is a critical design goal, it has been proposed to replace the steel vessels with composite vessel that weighs substantially less. A problem with using composite vessels is that it is difficult to make such vessels with lugs or other integral attachment devices.
Accordingly, a need exists for a mounting assembly for mounting a composite bottle or pressure vessel to a structure. The assembly must be able to accommodate the expansion and contraction of the vessel as it fills and empties.


A CR oscillation circuit comprises first and second logic elements, a capacitive element, and a resistive element. The first logic element is connected between a first node and a second node. The second logic element is connected between the second node and a third node. The capacitive element is connected between the first node and the second node. The resistive element is connected between the first node and the third node. The capacitive element includes a well, a diffusion layer, a gate electrode and a gate oxide film. The capacitive element has a voltage-dependence characteristic such that its capacitance value varies according to the variation in a supply voltage. The capacitance value of the capacitive element decreases when the on-resistance of the first and second logic elements increases according to the variation in the supply voltage.

1. Field of the Invention
The present invention relates to a CR oscillation circuit which generates a clock having a frequency determined according to the resistance value and capacitance value.
2. Description of the Prior Art
FIG. 12 shows a structure of a conventional CR oscillation circuit. The CR oscillation circuit includes inverters 91, 92 and 93, a capacitive element 94, and a resistive element 95. This circuit outputs clock CLK which has a frequency determined according to a product of the capacitance value of the capacitive element 94 and the resistance value of the resistive element 95 (see, for example, Japanese Laid-Open Patent Publication No. 5-206731 and Japanese Laid-Open Patent Publication No. 7-131301).
Clock CLK generated by the CR oscillation circuit is used as, for example, a display control clock for driving display elements provided over a display panel, such as a liquid crystal panel, or the like. It has been generally known that, if the frequency of the display control clock is varied, the timing of driving the display elements is also varied, so that flickers are seen in a representation of the display panel.
However, to keep the frequency of clock CLK constant as much as possible, the capacitive element 94 is desirably formed using a material which exhibits only a small variation in capacitance value even in the case of a variation in the supply voltage (a material having a small voltage-dependence characteristic). For example, the interpolysilicon (polysilicon-to-polysilicon) capacitance or intermetal (metal-to-metal) capacitance is used as the capacitive element 94. The resistive element 95 is desirably formed using a material which exhibits only a small variation in resistance value even in the case of a variation in the supply voltage (a material having a small voltage-dependence characteristic) and which exhibits only a small variation in resistance value even in the case of a variation in temperature (a material having a small temperature-dependence characteristic). For example, a polysilicon resistance is used as the resistive element 95.
However, since the polysilicon resistance has a small temperature-dependence characteristic, the frequency of clock CLK varies according to a variation in the operating temperature or environmental temperature of a semiconductor integrated circuit device incorporating the CR oscillation circuit.
In the case of using the interpolysilicon capacitance or intermetal capacitance, an additional step for forming the capacitance in a semiconductor integrated circuit (an exclusive polysilicon formation step for forming a second layer capacitance or an exclusive metal formation step for forming an intermetal capacitance) is necessary. Thus, the cost of the semiconductor integrated circuit increases.
Transistors included in the inverters 91, 92 and 93 also have the voltage-dependence characteristic, and therefore, the frequency of clock CLK varies due to a variation in the supply voltage. As a result, flickers occur in a representation of the display.

Matsushita Electric Industrial Co., Ltd.


A filtration and air supply device(3) used in an aquatic equipment(2) ,comprising filtration system with air supply, including filtration box(31), water pump(33), filter cartridge(32), air pump(7), air inlet tube(8), air duct(36). The characteristics of this filtration and air supply device(3) are: Said filtration and air supply device(3) locates inside the aquatic equipment(2). Box wall(310) is with grids(312). Water pump(33) locates in one side of the filtration box(31), and filter cartridge(32) locates in the other side of the filtration box(31). Bottom slab(9) is equipped with air inlet hole(38) and water pump power cord inlet hole(40). The bottom of air inlet hole(38) connects the air pump(7) through air inlet tube(8). The upper part of air inlet hole(38) connects air outlet hole (318) through air duct(36). No-return valve(37) locates between air duct(36) and air outlet hole(318). Outside of the air outlet hole(318) is equipped with air filter hood(315).

This invention involves a device used in aquatic equipment which contains a filtration and air supply system, especially combines water flow filtration and air supply into one system, and locates inside the aquatic equipment.
Aquatic equipment is used for keeping aquatic animals, particularly decorative fish, normally including water container, filtration system and air supply devices. The filtration system and air supply system of the existing aquatic equipment are separated located on the upper/side/outside of water container, which are in large volume, higher cost, can not be concealed and decorated. Electrical appliances, power cord and air tube are normally located inside or outside the water container, which made it look tangling and not nice, hard to clean and replace.

Filtration and air supply device used in an aquatic equipment


The invention relates to a spring core for mattresses and seat cushions with leaf springs. The spring core consists of a plurality of spring elements bent up out of passages punched in a plate, with each spring element having an upper and a lower spring clip, with the upper spring clip bent up out of a passage punched in the lower spring clip. The spring elements can also be coated with extruded foam.

1. Field of the Invention
The invention relates to a spring core for mattresses and seat cushions with leaf springs.
2. Description of the Related Technology
Spring cores for mattresses and seat cushions with spiral springs are known. These spring cores, which come in various embodiments, are not corrosion-resistant. They cannot be made of plastic but instead must be made only of spring steel wire.
It is also known that foam inserts and various natural fiber inserts with elastic fillers can be used as a substitute for metal spring cores in mattresses and seat cushions. These inserts have the disadvantage of very poor air circulation with the related unfavorable moisture regulation which can develop fungus. The insert material can also suffer from fatigue and deformation. In addition, extensive lathing is necessary for ventilation.
As is apparent in Austrian Patent No. 399,994, spring cores are known with leaf springs consisting of two plates arranged opposite one another with spring elements bent out of them. This design can be used only to a limited extent because the connecting webs between the passages punched in the plates do not permit a point load and thus do not permit limited or local-only deformation.


A method of adjusting the size of an area to be deleted in a pen-based computing system is provided. An electronic pen may include an eraser end. The eraser end may be used to contact the display surface of a computer system. Application of pressure to the display through the eraser end of the electronic pen may result in the deletion of the text indicated by the position of the electronic pen. As additional pressure is applied, the size of the area to be deleted increases. The invention provided may be used with both pixel type and stroke type erasing systems and provides an efficient and accurate method of deleting electronic ink.

Typical computer systems, especially computer systems using graphical user interfaces (GUIs), are optimized for accepting user input from one or more discrete input devices, such as a keyboard for entering text, and a pointing device, such as a mouse with one or more buttons, for operating the user interface. An example of such a GUI is the user interface for the WINDOWS® computer operating system (available from Microsoft Corporation of Redmond, Wash.). The ubiquitous keyboard and mouse interface provides for fast creation and modification of documents, spreadsheets, database fields, drawings, photos, webpages, and the like.
Recently, however, pen-based computing systems, such as tablet PCs, personal digital assistants, and the like, have been increasing in popularity. In pen-based computing systems, user input advantageously may be introduced using an electronic "pen" or stylus (e.g., akin to writing with a pen or pencil on a piece of paper). Indeed, in at least some pen-based computing systems, all user input is capable of being entered and manipulated using an electronic pen input device, and the user interface is fully controllable using only the electronic pen.
As pen based computing systems become more popular, users are increasingly entering more data in the form of electronic ink. As with any computer input device, errors are possible and corrections must be made. The ability to erase electronic ink from the user interface is necessary and the ability to use the electronic pen to erase this ink can be highly efficient.


An isolated switched-mode power converter converts power from an input source into power for an output load. Power switches within a primary-side power stage control the amount of power input to the power converter and, ultimately, provided to the output load. A digital controller on the secondary side of the power converter generates signals to control the power switches. This controller also senses a rectified voltage on the secondary side of the power converter and uses this sensed voltage to detect fault conditions of the primary side. For example, the sensed rectified voltage is used to detect undervoltage or overvoltage conditions of the input power source of the power converter, or faulty power switches within the primary-side power stage.

Isolated switched-mode direct-current (DC) to DC power converters use a transformer to convert power from an input source into power for an output load. Such power converters include power switches that convert DC input power into alternating current (AC) power that is fed to the primary side of the transformer. AC power supplied on the secondary side of the transformer is rectified and filtered so as to provide DC power to the output load. The primary-side power switches are typically controlled by pulse-width-modulated (PWM) control signals. A controller generates the PWM control signals with a frequency and duty cycle that are appropriate to meet the power needs of the output load.
The controller typically uses a linear closed-loop feedback technique to maintain the output voltage near a desired target. The controller may be implemented using analog or digital circuitry, and may be located on the primary or secondary side of the power converter. So as to maintain the integrity of the isolation barrier of the power converter, any signals crossing between the primary and secondary sides must pass through isolators, e.g., transformers, opto-couplers. Locating a controller, analog or digital, on the primary side of the power converter provides the advantage that the input voltage and/or current may be readily sensed and used to detect primary-side fault conditions. However, such primary-side controllers have the disadvantage that one or more analog isolators must be used to transfer the output voltage and/or current information from the secondary to the primary side for purposes of implementing the closed-loop control. Furthermore, system managers are typically located on the secondary side, meaning that communication signals between a primary-side controller and any such system manager must also pass through isolators. For these reasons, as well as cost, size, and flexibility advantages, isolated switched-mode power converters increasingly use digital controllers that are located on the secondary side. Detection of fault conditions on the primary side of a power converter presents difficulties for a controller located on the secondary side of the power converter.
Techniques are desired for detecting primary-side fault conditions in an isolated switched-mode power converter, wherein the techniques are readily implemented within a secondary-side controller. These techniques should require minimal or no additional circuitry relative to the circuitry required for typical operation of the power converter.

Voltage and current protection in isolated switched-mode power converters with secondary-side rectified voltage sensing


An automated chemical storage and feed system for a use station to chemically treat and monitor a process at a use station where the process is being run, including one or more base tanks at the use station storing chemicals and having means intermittently or continuously feeding the chemicals to the process. A level sensor is provided in each base tank to determine the chemical level and to deliver a signal to a control and processing unit. Sensors may be provided on the process for generating chemical condition information to the control and processing unit. A telemetry system receives the information from the control and processing unit and transmits it to a supply station where the information is recorded and read out by a person to monitor the chemical supply in the base tank and to monitor the chemical condition of the process so that the chemical feed may be adjusted as needed from the supply station and so that futher chemical may be delivered to the use station to be introduced into the base tank for maintaining an adequate supply to prevent interruption of the process. Refill tanks are shipped from the supply station to the use station, transported to the base tank, and emptied into the base tank. Thereafter, the refill tanks are returned to the supply station to be refilled. A unique hand truck is used at the use station to transport a refill tank to the base tank and to handle the refill tank during unloading.

Processes run by users/customers who require chemical additives/product must have a chemical supply stored at the process site where the process is being run and means for feeding the chemical into the process. For example, water-cooling towers that depend upon water for thermal conduction require chemical additives in the water in order to minimize corrosion of the equipment handling the water. Depending upon the chemical condition of the water, one or more chemicals may need to be added continuously or intermittently in order to produce the best possible results.
Heretofore, it has been customary to merely use a base supply tank at the use station for storing chemical and to intermittently or continuously feed the chemical into the water or liquid stream of the process, and provide one or more auxiliary refill tanks maintained on hand for refilling the base tank. The refill tanks, when empty, are then exchanged for other refill tanks shipped to the use station. Thus, an extra supply must always be maintained in inventory at the use station. Systems are used on site to determine the chemical condition of the water so that the proper feed of chemicals from one or more base tanks can provide the best treatment of the water for the process. When the base supply tanks are empty, they have been previously refilled if they are determined to be empty at the appropriate time either by the process workers or by personnel from the chemical supplier. Reliance on either has not been reliable, which then places the process in jeopardy and is costly to the user since the equipment condition deteriorates faster and the efficiency of the process is decreased. Further, there is always the chance of human error creeping into the ordering process for replenishment of the chemical supply. When improper control of the feed of chemicals to the process is encountered, such can cause use of excess chemicals, which is more costly to the user. Carrying an inventory of chemicals is also costly.

Automated chemical storage and chemical feed system


The present technology relates to a nanoparticle platform based on the unique and varied properties of DNA. Circular DNA can be replicated using a strand displacing polymerase to generate long linear concatamers of controllable length that spontaneously fold into a ball conformation due to internal base-pairing. These balls of DNA are discreet particles that can be made in variable sizes on a nanometer size scale in a scalable manner. The particles can be used in a variety of manners, discussed herein, including specific targeting, drug delivery to cancer cells, and diagnostics. Nanoparticles may also serve as multifunctional platforms for the integration of many currently used cancer therapeutic techniques.

DNA is one of the most thoroughly characterized molecules with regard to structure, chemistry, and modification, and has the capability to serve a wide variety of functions. It can be, and has been, used as a scaffold for the integration of varying entities due to its well defined ability to base-pair hybridize. With the discovery of aptamers and the realization that DNA has the capability to structurally and chemically recognize other molecules with near anti-body specificity at a fraction of the difficultly of synthesis, a new field of DNA based targeting molecules was born. Today there are FDA approved therapeutics based solely on DNA, such as the VEGF aptamer.
Current nanoparticle-based approaches to treating cancer include constructs composed of polymer, silica, or gold nanoparticles, liposomes, and less frequently such platforms as carbon nanotubes and viral capsids. These structures are coated with a variety of functionalizing entities such as polyethylene glycol (PEG) for biocompatibility, various targeting peptides, antibodies, small molecules, or aptamers and some form of therapeutic.

Single molecule nucleic acid nanoparticles

The Regents of the University of California


According to one embodiment, a desalination treatment membrane includes a desalting membrane and a base material which is disposed in close contact with the desalting membrane, wherein a solid salt is fixed to the base material by a graft-polymerization.

Reverse osmosis membrane (RO membrane) methods have hitherto been widely used in a desalination method of seawater. A reverse osmosis desalination method (RO method) is a method in which a pressure of about 55 atmospheres is applied to an osmosis membrane in a direction opposite to the osmotic pressure, thereby taking out fresh water from about 3.5% by weight seawater.
It is known that when a polymer electrolyte membrane obtained by a direct graft-polymerization of an electrolyte monomer to a hydrophobic polymer membrane is utilized in the RO method, a transmembrane flow rate is increased.
On the other hand, a forward osmosis membrane seawater desalination method (FO method) is known as the desalination method. According to this method, an osmosis membrane, which is the same as that used in the RO method, is used, and an aqueous ammonium carbonate having a higher concentration than that of seawater is disposed at a support membrane side, whereby fresh water is drawn into the aqueous ammonium carbonate side due to an osmotic pressure caused by ammonium carbonate, without applying a pressure. After that, the temperature of the ammonium carbonate solution is elevated to about 60° C. by heating it to decompose it into carbonic acid and ammonia, from which water is removed, thus resulting in acquisition of fresh water.
Further, there is also a composite semipermeable membrane having an ionic group or non-ionic group.


An on-chip electromagnetic (EM) pulse protection circuit detects EM pulse attacks, generates an alarm, and performs a defensive action to protect the integrated circuit. The EM pulse protection circuit can be used with various integrated circuits or manufactured chips in which, for example, there is a desire to keep information secure, maintain the security of the chip, secure boot processes, and/or protect private keys.

Hardware security is becoming a real concern in all product categories. Chips are being subject to security attacks using power, clock and EM pulses, to extract information from the chip. In an EM pulse attack on an integrated chip, attackers can use a probe with an inductor coil attached at the end of it and apply a very high voltage ramp for a very short duration of time. This creates electromagnetic pulses at the tip of the inductor coil. Chip package as well as metal routes (VDD/GND rails) on a chip have some inductance associated with them. When such a pulse is applied and directed to a chip via the probe, the applied pulse can produce glitches on the chip due to the mutual coupling between the probe coil and the on-chip inductance. The glitches can be caused even without the probe directly touching the chip and when the mutual coupling is small. The induced glitches can manifest into timing failure if the glitch is big enough or a clock glitch which can cause a unit to skip a cycle or corrupt data that allows attackers to extract secure information.
If an EM pulse attack is not successful at one location, attackers can keep sweeping the probe for target locations on the chip and/or increase the attack voltage to successfully break the chip. The attackers can use a robotic arm having an XYZ-positioning mechanism in order to target a particular location on the chip.

Protecting chips from electromagnetic pulse attacks using an antenna


Illumination device with at least one LED as light source, the LED emitting primary radiation in the range from 380 to 420 nm, this radiation being partially or completely converted into longer-wave radiation by phosphors which are exposed to the primary radiation from the LED, the conversion taking place at least with the aid of a phosphor which emits over a broad band and originates from the class of the (Eu,Mn)-coactivated halophosphates, the cation M representing one of the metals Sr, Ca, Ba.

The invention is based on an illumination device with at least one LED as light source in accordance with the preamble of claim 1. It deals in particular with an LED which emits in the visible or white region and is based on an LED which emits primarily UV/blue light.
An illumination device with at least one LED as light source which emits, for example, white light is currently produced predominantly by combining a Ga(In)N LED which emits in the blue region at approximately 460 nm and a yellow-emitting YAG:Ce3+ phosphor (U.S. Pat. Ser. No. 5,998,925 and WO 98/12757).
A more complex concept with improved color rendering is a three-color mixture. In this case, the fundamental colors red-green-blue (RGB) are mixed to produce white. In this context, it is either possible to use a blue LED for partial conversion of two phosphors which emit red and green (WO 00/33390) or to use a UV-emitting LED which excites three phosphors, which respectively emit in the red, green and blue (cf. WO 97/48138). Examples are line emitters, such as YOB:Ce,Tb (green) and YOS:Eu (red). In this case, however, a relatively short-wave emission (UV region less than 370 nm) is required in order to be able to achieve high quantum efficiencies. This requires the use of sapphire substrates for the UV LED, which are very expensive. On the other hand, if a UV LED based on the less expensive SiC subtrates is used, one has to be satisfied with an emission in the range from 380 to 420 nm, which makes the use of line emitters in the green and red difficult to impossible. In the case of blue phosphors, this leads to absorption problems.
Furthermore, a specific problem in this context is the additional absorption loss of blue radiation on account of the broad-banded nature of the absorption of the red- and green-emitting phosphors. All this together leads to considerable limitations when setting the luminous color and light efficiency.
U.S. Pat. No. 5,535,230 has disclosed the use of ordinary xe2x80x9ccalcium halophosphatexe2x80x9d for semiconductor components or fluorescent lamps. This term is understood as meaning the commercially available Ca5(PO4)3Cl:(Sb,Mn), as discussed in detail, for example, in U.S. Pat. No. 5,838,101. A further halophosphate discussed in that document is (Sr,Ba,Ca)5(PO4)3Cl:Eu. In this context, these phosphors are used for excitation by the far UV (254 nm).
Finally, U.S. Pat. No. 3,513,103 proposes the use of (Sr,Ca)9(PO4)6Cl2:(Eu,Mn) in fluorescent lamps for conversion of the UV radiation.
It is an object of the present invention to provide an illumination device with at least one LED as light source, the LED emitting primary radiation in the range from 380 to 420 nm, this radiation being partially or completely converted into longer-wave radiation by phosphors which are exposed to the primary radiation from the LED, which is distinguished by high stability. A further object is to provide an illumination device which emits white light and in particular has a high color rendering. A further object is to provide a high-efficiency LED which absorbs well in the range from 380 to 420 nm and is easy to produce.
This object is achieved by the following features: the conversion is achieved at least with the aid of a phosphor which emits over a broad band and originates from the class of the (Eu,Mn)-coactivated halophosphates, the halophosphate corresponding to the formula M5(PO4)3(Cl,F):(Eu2+,Mn2+), where M=Sr, Ca, Ba individually or in combination. Particularly advantageous configurations are given in the dependent claims.
High-efficiency phosphors which emit red light and are suitable for LEDs which emit primarily in the near UV are rare. The term near UV is in this context intended to cover wavelengths between approximately 380 and 420 nm. By way of example, the line emitter Ln2O2S:Eu3+ (Ln=La, Y, Gd, Lu) is suitable for use with UV LEDs with primary emission below 380 nm. Hitherto, only Eu-doped nitrides or sulfides have been serious candidates for wavelengths above 380 nm. However, nitrides, which are sensitive to air, are difficult to produce. Sulfides have two drawbacks: they have an absorption gap between 380 and 410 nm and they are not particularly stable, i.e. are temperature-sensitive. Organic phosphors are too sensitive for LED applications. Virtually all known oxides have a low absorption at wavelengths greater than 380 nm.
According to the invention, the phosphor used for the LED-based illumination device is a halophosphate of type (Sr,Ba,Ca)5(PO4)3(Cl, F):(Eu,Mn). The same halophosphate, but activated only with Eu, is well known as a blue-emitting phosphor for fluorescent lamps and is commercially available. It has very good absorption at 254 and 365 nm. However, its absorption in the near UV (380 to 420 nm) is unsatisfactory. Quite unexpectedly, it has been found that adding the coactivator Mn hugely increases the absorption in the desired region without the efficiency suffering significantly.
In particular, it has been found that adding a melting agent in a high concentration (SrCl2) and the use of a relatively high level of Eu are favorable. Furthermore, a high level of purity and a small grain size of the raw materials should be ensured. Then, depending on the concentration of the activators Eu and Mn, a yellow to red emission results. The Eu content based on the divalent cation is preferably 0.5 to 2 mol %, and the Mn content, based on the divalent cation, is 3 to 8 mol %. The cation used is primarily Sr, if appropriate with a small amount of added Ba and/or Ca. The Mn:Eu ratio should be approximately 4 to 8.
Overall, therefore, it is possible to achieve highly efficient conversion of primary radiation in the range from 380 to 420 nm, since this excitation range is close to the emission of the phosphor.
In addition to producing a colored light source by excitation by means of UV radiation or blue primary emission from an LED (380 to 420 nm), it is advantageous in particular to generate white light with the aid of this phosphor. In the case of a UV-emitting LED as primary light source, this is achieved using at least two phosphors. In this way, it is possible to achieve good color rendering of over Ra=75. The point is that the phosphor according to the invention itself simultaneously provides two emission peaks in the blue and the yellow-red spectral regions, and therefore achieves something which normally requires two different phosphors.
This phosphor can be used together with other phosphors to achieve a white-emitting LED. It is particularly preferable for it to be used together with chlorosilicate phosphors as are already known per se. In this context, see for example WO 01/95400, corresponding to U.S. Ser. No. 10/048,963 and WO 01/93342, corresponding to U.S. Ser. No. 10/031,578, to which reference is expressly made. Therefore, the white point can be hit almost exactly given a suitable mixture.
A white mixture can also be produced on the basis of a UV-emitting LED by means of this (Eu,Mn)-doped halophosphate together with a blue-green phosphor, such as for example BaMgAl10O17:Eu2+ (BAM). If necessary, the color rendering can be improved still further by adding a green phosphor (for example Eu-doped thiogallates or chlorosilicates or Sr aluminate).

Patent-Treuhand-Gesellschaft fxc3xcr Elektrische Glxc3xchlampen mbH


A spinal fixation device includes a device body configured for implantation in a human body and a clamp that is formed from shape memory material and attached to the device body.

Various embodiments of a spinal fixation device are described herein. In particular, the embodiments described herein relate to an improved spinal fixation device used to stabilize and fuse the human spine.
Devices and methods for correcting spinal deformities are known and use implants anchored to portions of the posterior spine at various locations on the vertebrae, such as on the lamina, transverse process, or pedicle, as shown in FIG. 26. Such known devices stabilize the spine by fusing the spinal column into a single rigid construct, and typically include pedicle screws or hooks that are anchored to the spine and connected by fusion rods.
The placement of pedicle screws requires a large amount of time and great effort by a surgeon, and placement location options are limited. A surgeon must first very exactly position and form a screw hole in the pedicle bone. The hole must be tapped and then the pedicel screw must be screwed into the tapped hole in the pedicle bone.
Known spinal hooks may be placed at various locations of the vertebrae, including on the lamina, transverse process, and pedicle. The design of known spinal hooks requires the surgeon to place two hooks in an opposing relationship, and then tighten each hook to the fusion rod to hold the hooks to the vertebrae. The size and shape of the portion of the spinal hook that engages the vertebrae does not perfectly mate or engage with the portion of a vertebra to which it will be attached. Further, known spinal hooks have no stabilization mechanism and are therefore not held to the vertebrae until a fusion rod is attached between two or more spinal hooks. Known spinal hooks can therefore be unstable during surgery prior to and during attachment of the fusion rod, and may fall off the bone during surgery prior to the fusion rod being attached, causing difficulty for the surgeon. Accordingly, there is a need for more stable and reliable spinal hooks.


A liner for lining a receptacle has a front panel and a rear panel that define an interior and a top opening in fluid communication with the interior. A frangible sizing seal structure joins together the front and rear panels along the top of the liner adjacent a side of the liner. The sizing seal structure limits a size of the liner opening and is configured to be selectively broken to separate the front and rear panels to enlarge the size of the liner opening so that the top edge margin fits over a rim of the receptacle. In a method of making the liner, the front and rear panels are sealed together to form the frangible sizing seal structure.

Liners are used to line waste cans and other receptacles. Conventional liners have closed bottom ends and open top ends. Typically, when liners are placed into a receptacle, the top end portion of the liner is loosely fitted around the rim to secure the liner to the receptacle. Conventional liners lack features for securely gripping the rim of the receptacle and are, therefore, prone to fall off in use.


An adaptive control system for an electric motor has energization circuitry connected to phase windings of the motor for energization thereof and a controller for generating a control signal corresponding to a profile of an excitation current. The control signal is applied to the energization circuitry to control energization of the phase windings. The energization circuitry provides the excitation current to the phase windings from a power supply. For a present combination of torque and speed, the controller adaptively determine the profile of the excitation current optimal for achieving a particular motor control objective. Excitation current profiles may be optimized to achieve maximum efficiency, maximum torque, minimum torque ripple, minimum core loss, etc.

Depending on motor's applications and user's needs, a motor control strategy may have different objectives. Motors may be controlled to achieve maximum torque, maximum torque per ampere, minimum core loss, maximum efficiency for a given combination of torque and speed, minimum torque ripple for a given combination of torque and speed, maximum constant power speed ratio (CPSR), or other objectives.
For example, recent advances in high-energy batteries, combined with the development of smaller and more powerful motors for electric vehicles, portable appliances and electronic equipment, emphasize the need for a high-efficiency motor control scheme to achieve low power consumption in order to extend battery life, providing longer hours of operation. However, in some applications, such as vehicle driving operation, there may be a need for torque capability in excess of that available from the most efficient motor control scheme. A driver may demand higher torque than the system can accommodate with the high-efficiency motor control scheme. In servo control applications, or when smooth control is required at lower speeds, minimizing the torque ripple becomes the main issue for a control strategy. For control of motors that require wide ranges of constant power operations, such as adjustable-speed driving motors, it is necessary to achieve high constant power speed ratios in the vicinity of 10 or even higher.
Motor control parameters may be optimized to achieve desired objectives. However, variable conditions of motor operation and changing motor's parameters call for frequent re-evaluations of motor control parameters optimized to achieve certain objectives. The need thus exists for a motor control system able to adaptively optimize motor control parameters during motor operation.

Adaptive system for optimizing excitation current waveform profiles for electric motors


A method, computer program product, and computing system for dividing a plurality of volumes replicated across a pair of storage systems into one or more consistency groups. A polarization state may be defined for each consistency group. An input-output (IO) failure associated with at least one consistency group may be detected. At least a portion of the at least one consistency group may be polarized based upon, at least in part, the polarization state defined for the at least one consistency group.

Storing and safeguarding electronic content may be beneficial in modern business and elsewhere. Accordingly, various methodologies may be employed to protect and distribute such electronic content.
For example, multiple storage systems may be configured in an active/active relationship where the two storage systems present a single volume. The volume data may be available through paths going to both storage systems. Every write operation to one storage system may be replicated to the other storage system synchronously. The two storage systems may maintain a cluster between them (e.g., via a Time to Live (TTL) mechanism, where a "preferred" storage system grants TTL to its peer, the "non-preferred" storage system). In this configuration, once the TTL expires for the non-preferred storage system, the non-preferred storage system may go offline. In particular, once it is impossible to replicate IO write requests between the two storage systems within the TTL, one of the storage systems must stop servicing all IO requests to the host before allowing the other storage system to proceed. This process is generally referred to as polarization.
However, when multiple volumes are replicated between storage systems, problems may arise when polarizing the volumes. For example, polarizing individual volumes may result in data loss for groups of volumes, while polarizing large portions or even all volumes may result in significant recovery times and performance penalties for the storage systems.

System and method for volume polarization across multiple storage systems


An aircraft de-icing system and method for removing the ice from the surface of an aircraft. The aircraft de-icing system reduces power requirements by allocating power from a single power source between the parting strip elements and the shed zone elements. The aircraft de-icing system may utilize a measurement of heat transfer to schedule the time necessary to shed ice from the surface of the aircraft. The de-icing system may also utilize a method for calculating the ice accumulation on the surface of the aircraft based on the outside air temperature.

1. Field of the Invention
The invention relates generally to aircraft. More specifically, the invention relates to the field of de-icing systems and methods for aircraft.
2. Description of the Related Art
Traditional electrical de-icing systems for aircraft employ parting strips and shed zones to prevent ice buildup on various elements of the aircraft during operation. The parting strips, on the leading edge of aircraft surfaces, and shed zones, on the upper surface of the aircraft, incorporate heating elements to heat the surface of the aircraft. This causes ice deposited on the surface to be de-bonded and be removed by aerodynamic forces. The use by the de-icing system of both parting strips and shed zones may lead to excessive power requirements to operate both types of elements concurrently. Such de-icing systems also utilize a number of sensors to determine the environmental conditions in which the aircraft is operating and the state of various elements of the de-icing system increasing the complexity of the de-icing system and its mean time to failure.


A write-behind enablement approach is presented which allows writing data to multiple volumes of storage media associated with one or more server nodes in a distributed processing environment. A client application on a client node writes blocks of data to a first magnetic tape of a first server node until an end of physical storage of the first magnetic tape is reached, without having predetermined the size of the first magnetic storage. Thereafter, the writing of blocks of data is switched to a second magnetic tape, which may be on the first server node or a second server node in the system. The writing and switching process can repeat a number of times across multiple magnetic tapes. Data is written in a write-behind operation and the switching to a new magnetic tape occurs transparent to the client application, and without loss of data. The client and server nodes are adapted to ascertain how many blocks of data were written to the first storage medium, and the second storage medium, and to allow for the client process to write a header label to the second storage medium prior to writing blocks of data thereto.

In a distributed client/server storage environment, a client process will send data for storage to a remote server node, which writes the data to the actual storage medium. The ability to write data to remote storage medium allows an application program to use hardware that is associated with processors other than the one the application program is running on. The sharing of hardware resources in this way is a main advantage of a distributed computing environment. In order to improve performance within such an environment, it is desirable that the client program write several blocks of data to the server node for storage before reading a response from the server node to see if the transferred blocks of data have been written correctly. This technique of writing to the server node multiple blocks of data before confirming at a predetermined interval that the data has been correctly written is referred to as a "write-behind" function. With this function, communication overhead and delays between the client and server are minimized.
One difficulty with the write-behind process occurs because the client node does not continually check for status after sending each block of data. It will not be notified that the storage medium is full until the predetermined number of blocks have been sent to the server. Thus, the client node may only be notified of an end of volume condition after it has sent several additional blocks of data following the first block of data that would not fit in the storage medium.
One way that previous products have handled this problem is to require the client program to know the size of the storage medium before writing to it. As the client writes blocks of data to the storage medium, it knows that after a pre-specified number of blocks it must stop writing, and transition to a new storage medium. By doing this, the client never sends extra blocks to the server that cannot be written to the existing medium, and the problem is avoided.
There are two penalties for requiring the client to know the size of the storage media ahead of time. The first is that it is inconvenient to the programmer who must specify the storage medium size. The second is that in order to be safe, the maximum number of blocks is often not written to the storage medium. Thus, part of each storage medium may go unused.
In view of the above, a different approach to data transfer interfacing of the client and server nodes is desired.

Method for multi-volume, write-behind data storage in a distributed processing system

International Business Machines Corporation


A method and apparatus for forming runners (47, 56) in a mold assembly (10, 50). The mold assembly (10, 50) has a cavity plate (11, 51) and a sealing plate (12, 52). The sealing plate (12, 52) has a sealing plate cavity (18) in which a resilient material (14) is placed. A floating seal (13) is movably mounted to the sealing plate (12) such that the floating seal (13) is on the resilient material (14). A workpiece (19) is placed on the resilient material (14) and the cavity plate (11, 51) is mated with the floating seal (13) thereby forming the runners (47, 56). The resilient material (14) imparts a force on the floating seal (13) and pushes the floating seal (13) against the cavity plate (11, 51) to prevent leakage of encapsulating material from the runners (47, 56).

The present invention relates, in general, to mold assemblies, and more particularly, to a method and apparatus for forming a runner in a mold assembly.
Workpieces, such as semiconductor devices, are encapsulated to protect them from becoming damaged by mechanical stresses, moisture, or the like. In parting plane mold assemblies, encapsulation is accomplished by placing a workpiece in a first section of a mold assembly and mating the first section with a second section, wherein the step of mating encloses the workpiece in a mold cavity and forms a plurality of runners extending from the mold cavity. Subsequently, an encapsulation material, commonly referred to as a molding compound, is injected through the plurality of runners into the mold cavity, thereby encapsulating the workpiece.
Although this technique has been useful for encapsulating workpieces, it requires machining the mold assembly sections and forming the workpieces to exacting tolerances. A particularly critical problem in parting plane mold assemblies is manufacturing the mold sections so that there are no spaces or gaps along the mating surfaces. Spaces or gaps prevent proper sealing between the mold sections and allow seepage of the encapsulation material from runners formed when the mold sections are mated, i.e., the formation of flash. In addition, workpieces whose dimensional variations exceed specified limits may also cause flash formation by preventing the proper sealing of the runners. Since additional processing steps are required to remove the flash, dimensional variations of either the mold assembly or the workpieces increases the costs associated with encapsulating workpieces.
Accordingly, it would be advantageous to have a method and apparatus for forming a runner in a mold assembly that eliminates flash formation. The method and apparatus should create a seal between two mold sections by imparting a uniform pressure along a parting plane formed by mating the two mold sections. It would be of further advantage for the method and apparatus to prevent flash formation by compensating for dimensional variations of the components of the mold assembly or the workpieces.

Method and apparatus for forming a runner in a mold assembly


An electrical connector assembly having a housing and a cover assembly. The housing defines a cable receiving cavity and includes a conductive connector shield extending within a portion of the cavity. The cover assembly has a back cover that includes an end wall that mates with the housing adjacent to the cable receiving cavity; a cable support wall extending from the end wall and defining a cable passage extending through the end wall, with the cable support wall configured to receive a shielded cable therein; and a conductive spring clip fixed to the cable support wall and having a first leg configured to be biased into contact with a conductive ferrule of the shielded cable and a second leg, on an opposed side of the cable support wall, spring biased into contact with the connector shield.

This invention relates to an electrical connector, and more particularly to an electrical connector having electromagnetic interference (EMI) shielding.
Electrical connectors are known that provide shielding and grounding for electrical cables to which they connect. It is desirable to provide such electrical connectors with reliable shielding and grounding while minimizing the number of components to assemble/disassemble and allowing for ease of assembly and disassembly.

EMI grounding back cover for electrical connector


A grass and brush mowing attachment for a fork lift or other front loading machine which is operated hydraulically and which includes forwardly projecting tines. The mowing attachment has a deck with side skirts having runners on their lower edges. Sleeves on the deck receive the tines of the fork lift, and chains are also connected between the deck and the machine. The mowing deck has a blade which is driven by a reversible hydraulic motor that has detachable hose connections with the hydraulic pump on the fork lift. The blade is sharp on both edges and is equipped with pivotal kick back blades on its opposite ends.

This invention relates to mowing devices and, more particularly, to a hydraulically-powered mower that is attachable to the tines of front loader equipment, such as fork lifts or skid loaders.
One type of mower presently employed for cutting light brush and grass is designed to be connected by a 3 point hook up to the rear of a tractor. The Power Take Off (PTO) of the tractor is used as the source of power to rotate the cutting blades.
One drawback of this type of mower is the mower position. The mower is pulled by, and thus trails, the tractor. This results in a mower position which is not in the operator's preferred line of sight, which can be hazardous. Since it is often necessary to watch the mower to make sure it is performing adequately, the operator must turn around and, accordingly, momentarily look in a direction opposite from that in which the apparatus is traveling. This leaves the operator unaware of objects in the mower's path. Also, mowers commonly throw debris out the front portion of the mower, resulting in debris being thrown toward the operator.
Another drawback of this type of mower is the limited output efficiency of its power source in comparison with a hydraulic-powered device. Additionally, the operation of conventional mowers typically requires complex parts such as a gear unit and drive shaft. Each of these entails a variety of working parts which can be costly to replace in the event of a breakdown.
Another disadvantage of typical existing mowers is the inability of the blade to rotate in either direction. Additionally, a large area is required to turn the tractor and mowing device around.
The use of hydraulic power to drive the blades of a mower has been shown in U.S. Pat. No. 3,946,543 (Templeton). However, this device additionally employs chains, sprockets, and gears which are not as efficient or dependable as the coupling used in accordance with the present invention.
U.S. Pat. No. 4,487,006 (Scag) discloses a lawn mowing device having a cutting assembly attached by forks to a riding unit. The mower uses a gas engine to power its wheels for movement. The same engine is appropriately connected for rotating the cutting assembly. This mowing unit is designed as a whole, rather than an inexpensive mowing attachment.