An adjustment system and method is provided for determining a set position of a motor-driven adjustment part of a motor vehicle by a Hall sensor, which is arranged in a rotated magnetic field. The drive of the adjustment part is activated for a specified number of approaches of the adjustment part toward an end position. A counter is changed by a counting value upon each approach of the adjustment part toward the end position. A standardization drive to the end position is carried out when the specified number of approaches is reached, and each deactivation and/or activation of the Hall sensor is detected and the counter state of the counter is changed according to the number of deactivation and activation cycles.

1. Field of the Invention
The invention relates to a method for determining the set position of a motor-driven adjustment part of a motor vehicle by means of a Hall sensor, which is arranged in a magnetic field rotated as a result of a motor rotation, whereby the drive of the adjustment part is activated for a specified number of approaches of the adjustment part toward an end position, whereby a counter is changed by a count value upon each approach of the adjustment part toward the end position, and whereby a standardization run to the end position is carried out when the specified number of approaches is reached. It refers further to an adjustment system operating according to this method.
2. Description of the Background Art
A plurality of adjustment devices or adjustment parts (adjustment elements), driven by electric motors, are typically present in a modern motor vehicle. These are, for example, an electric window lift or a sunroof. During a setting process of such an adjustment part a desired end position is to be reached precisely, to which end the set position of the adjustment element must be known accurately. Moreover, it is also often necessary to know the current set position or the parameters that can be derived therefrom, such as the setting speed or the traveled adjustment path, for the reliable detection of a case of pinching.
For example, DE 199 16 400 C1, which corresponds to U.S. Pat. No. 6,225,770, and which is incorporated herein by reference, discloses the provision of a position and rotation sensor for the most precise detection possible of the set position of a window pane. It includes two Hall sensors offset at a distance or angle from one another or a Hall sensor (Hall IC) having two Hall probes (sensitive surfaces) and a multipole, for example, two- or four-pole ring magnet, which is arranged on the drive shaft of the electric motor. The Hall sensor detects a change in the magnetic field due to a rotation of the ring magnet, fixedly connected to the drive shaft, and generates counting pulses therefrom. These are evaluated together with information on the rotation direction of the ring magnet and thereby of the electric motor, by counting the counting pulses upward or downward depending on the rotation direction of the drive and thus indicating the respective position of the window pane.
DE 10 2006 043 839 A1 discloses a conversion of magnetic field changes, occurring at Hall probes or Hall sensors, of a comparator circuit with hysteresis (Schmitt trigger circuit) into two binary pulse trains, offset, for example, by 90° to one another. An upper switching threshold and a lower switching threshold are provided in such a comparator circuit with hysteresis. The motor or drive speed and thereby the position of the adjustment part and the rotation direction of the electric motor or rotary drive can be determined by counting the pulses per unit time by a comparison of the two pulse trains.
In adjustment systems with two end positions (open and closed), such as particularly in the case of a window lift drive, first as part of a first learning or standardization cycle, at least one of the end positions is directly approached and in this position the count value of a counter is set to a specific value, particularly to zero. During the movement toward the other end position due to an activation of the drive or electric motor in the opposite direction, the Hall sensor generates a number of counting pulses, which, for example, increment a position counter, whereby when the other end position is reached the current counter state is assigned to this end position and, for example, is stored. Thereby, each position of the driven adjustment part along the adjustment path between the two end positions is assigned a count value of the position counter. During the subsequent movements toward this end position, particularly the closed position (approach cycles), the drive is stopped shortly before reaching the standardized count value during normal operation in order to perform so-called soft starts of the adjustment element in the end stop, also designated as a capture region or seal entry, to increase the lifetime solely due to the inertia of the system.
To take into account changes in the adjustment system that occur during proper use, e.g., the end position, and to postadjust the adjustment system, DE 195 27 456 B4, which corresponds to U.S. Pat. No. 5,777,446, discloses the movement toward the end position in a standardization run without a soft start after a specified number of approach cycles and the resetting of the counter for the purpose of standardization to a specific value, for example, again to zero. In this way, a fixed number of approaches is specified after which a standardization run occurs in each case.
DE 198 38 293 A1, which corresponds to U.S. Pat. No. 6,346,787, discloses, moreover, instead of a fixedly specified number of approach cycles, to design the number of cycles as variable and to this end after each standardization run to generate a random number for the subsequent approach cycles, after which the next standardization run is to occur. The wear of the system is said to be reduced and the lifetime increased by generating the random number, which is decremented with each approach cycle of the adjustment part, to be positioned, up to the specified count value, which triggers a standardization run.

Method for determining the set position of an adjustment part

Brose Fahrzeugteile GmbH & Co. KG, Hallstadt


A coin selector has a coin passageway wherein coins are inserted. A magnetic sensor inspects the coins in the coin passageway, and generates an inspection signal, which a controller evaluates to generate an acceptance or rejection signal. A gate plate responds to the acceptance or rejection signal, and selectively provides access to one of two chutes for the inspected coins, namely, an accepting chute for acceptable coins and a return chute for unacceptable coins. A photo sensor detects the acceptable coins in the accepting chute, and generates a detection signal. Counters are adapted to count the acceptance signals incrementally and the detection signals decrementally. A controller adds together the two counts of the counters, evaluates the sum of this addition, and generates an error signal in accordance with the sum evaluation. The error signal is utilized for the detection of fraud.

1. Field of the Invention
The present invention relates to a coin selector and selection method for coin-operated machines for detecting fraud in coin insertion. More particularly, the present invention relates to a coin selector applied to a coin-operated machine such as a gaming machine, and relates also to an error detecting method to detect fraud in coin insertion through the same selector.
2. Description of the Prior Art
A coin-operated machine is operated in response to insertion of coins, tokens, medals or other disks (herein referred to as coins) into a coin inlet or passageway. Such coin-operated machines include slot machines, other gaming machines, vending machines and money-changing machines. In a slot machine for example, an acceptable coin must be selected before a game can begin, in order to prevent the machine from operating by insertion of a slug or unacceptable coin different from the genuine acceptable coin of a predetermined denomination.
To automate the selection of acceptable coins for the slot machine, there has been proposed a coin selector as illustrated in FIG. 6. In this conventional coin selector, a coin 9 is sensed through an inlet slot of the slot machine and passed along a coin passageway 10, which communicates both with an accepting chute 11 and a return chute 12. The coin 9 is detected twice: by a magnetic sensor 13 and by a photo sensor 14. A gate 15 is swingable between a position wherein the coin 9 from the passageway 10 is passed to the accepting chute 11, shown in phantom line in FIG. 6, and another position wherein the coin 9 is passed to the return chute 12 shown in full line in FIG. 6.
An inspection signal of the coin 9 is sent from the magnetic sensor 13 to a controller 16, and is evaluated as to whether it represents an acceptable or unacceptable coin. When the coin 9 is acceptable, the plate 15 is moved to the accepting position by actuation of a solenoid 17, so as to guide the acceptable coin 9 into the accepting chute 11. The acceptable coin 9 then passes the photo sensor 14 to cause it to output a detection signal to the controller 16, which evaluates the detection signal as to effectiveness. To conduct the latter evaluation, the controller 16 judges a time period T1 which begins with inspection signal and ends with the detection signal. When the detection signal is judged as effective, the controller 16 actuates a slot machine mechanism 18. The player can play games corresponding to the number of coins inserters.
The conventional coin selector, however, suffers from the disadvantage of being vulnerable to fraud by use of a tool as illustrated in FIG. 7, enabling a player to play games without paying coins. A celluloid plate 19 of the tool is crooked in correspondence with the inside of the passageway 10, and is provided with three slots 19a to 19c in positions corresponding to the downstream sensor 14. An acceptable coin 9 as a decoy is secured to the celluloid plate 19 in the position corresponding to the upstream sensor 13. The coin 9 is evaluated as acceptable by the sensor 13 and the controller 16. The slots 19a to 19c are detected by the photo sensor 14, and cause the controller 16 to generate three detection signals before the lapse of time T1 from the sensing of the coin 9 by the sensor 13, and are evaluated as effective three times. One use of the celluloid plate 19 thus enables the user to play three games without actually paying any coins.

Coin selector and selection method for coin-operated machines for detecting fraud in coin insertion


For use in a mobile ad hoc network (MANET), a first MANET node for aggregating route cost information associated with a first route from a source MANET node to a destination MANET node. The first MANET node comprises: 1) a radio frequency (RF) transceiver for wirelessly communicating with other MANET nodes; and 2) a controller for receiving incoming data packets from the RF transceiver and sending outgoing data packets to the RF transceiver. The controller receives a Route Request (RREQ) message generated by the source MANET node and retrieves initial route cost information from the RREQ message. The initial route cost information comprises an RF link cost parameter for a preceding RF link between the first MANET node and the source MANET node and a node cost parameter for a preceding MANET node between the first MANET node and the source MANET node.

Wireless network topologies that enable wireless nodes (i.e., mobile stations, wireless terminals) to communicate with each other and with fixed networks generally fall into two categories: i) infrastructure-based and ii) infrastructureless. Infrastructure-based networks have traditionally been based on the cellular concept and require a high level of infrastructure support. In an infrastructure-based network, wireless nodes communicate through access points (e.g., base stations) connected to the fixed network (e.g., Internet). Typical infrastructure-based networks include GSM networks, UMTS networks, CDMA networks, WLL networks, WLAN, and the like.
In an infrastructureless network, wireless nodes (i.e., mobile stations, wireless terminals) communicate directly rather than through access points or other base stations. One common and increasingly popular infrastructureless network topology is the mobile ad hoc network (MANET). A MANET is a group of wireless nodes that dynamically form a network with each other and without using any pre-existing fixed network infrastructure. In most cases, wireless nodes of a MANET are small mobile devices that are relatively limited in term of CPU capability, memory size, and power consumption.
Mobile ad hoc networks are expected to continue to grow over the next 2-3 years to become the dominant peer-to-peer communication technology. Cell phones are being equipped with IEEE-802.11 and other wireless LAN technologies. The proliferation of cell phones and the ubiquity of economical IEEE-802.11 networks will create a new kind of mobile, ad-hoc and peer-to-peer network.
As noted above, a mobile ad hoc network (MANET) can be set up without using a pre-existing network infrastructure. This can be done anywhere and at anytime. The wireless nodes of a MANET are connected by wireless links and are free to move randomly. The wireless nodes also act as routers. A MANET supports traffic types that are different from the traffic types typically found in an infrastructure-based wireless network. MANET traffic types include: 1) peer-to-peer traffic; 2) remote-to-remote traffic; and 3) dynamic traffic.
In peer-to-peer traffic, there is only one hop between the communicating wireless nodes (i.e., direct communication). In this instance, the network traffic (in bits/second) is usually constant. In remote-to-remote traffic, there are two or more hops between communicating wireless nodes, but a stable route is maintained between the source and destination nodes. This often occurs if several nodes stay within range of each other in one area or if the nodes move as a group. Dynamic traffic results when the MANET nodes move around and communication routes must be reconstructed. This often causes poor connectivity and network traffic occurs in short bursts.
Each MANET node is autonomous and may function as both a host and a router. Thus, each wireless node performs basic host processing and performs router switch functions. Thus, endpoints and switches are indistinguishable in a MANET. Since there is no central network to control network operations, control and management of a MANET is distributed among the wireless nodes. The MANET nodes cooperate to implement security and routing functions.
A MANET may implement different types of routing. Basic types of ad hoc routing algorithms are single-hop and multi-hop. These are based on different link layer attributes and routing protocols. A single-hop MANET is simpler than a multi-hop MANET, but lacks the functionality and flexibility of a multi-hip MANET. When delivering data packets from a source to its destination out of the direct wireless transmission range, the packets should be forwarded via one or more intermediate nodes.
Since MANET nodes are mobile, the radio frequency (RF) links may change rapidly and unpredictably over time. In order to compensate for traffic and propagation conditions, the MANET nodes dynamically modify routing information between each other as the nodes move, thereby forming new network topologies. Advantageously, a MANET node may operate not only within the mobile ad hoc network, but also may access a public fixed network (e.g., cellular network).
MANET nodes may use, for example, an Ad Hoc On-Demand Distance Vector (AODV) routing protocol in ad hoc network environments in which the movement of the MANET nodes cause frequent changes in RF link quality. The AODV protocol enables the MANET nodes to adapt quickly to dynamic link conditions. The AODV algorithm enables dynamic, self-starting, multi-hop routing between mobile MANET nodes in an ad-hoc network. The AODV protocol enables the mobile MANET nodes to obtain routes for new destinations quickly and does not require the MANET nodes to maintain routes to destinations that are not in active communication. The AODV protocol provides quick convergence when the ad hoc network topology changes (e.g., a new MANET node joins the network).
The AODV protocol uses a destination sequence number for each route entry. The destination node creates a destination sequence number for any usable route information the destination node sends to a requesting node. Using destination sequence numbers ensures loop freedom (i.e., prevents loops). Given a choice between two routes to a destination MANET node, a requesting node always selects one with the greatest sequence number. Advantageously, when the AODV protocol detects an RF link breakage, the AODV protocol immediately transmits notifications only to the affected set of nodes. Other nodes are not notified.
The current state of the art in MANET routing algorithms can be classified into two broad types: i) pro-active protocols and ii) reactive protocols. A proactive protocol typically uses a single aggregate value as the route cost between a source MANET node and a destination MANET node. Unfortunately, this means that fine-grained information is lost. Furthermore, the routing table in a pro-active protocol contains only one best route that is associated with the single aggregate value. This is a disadvantage because two different applications on the same node might have different quality-of-service (QoS) requirements. For example, an FTP application requires large bandwidth and a low error rate and tolerates high latency. However, a voice application requires low bandwidth and a low latency and tolerates a moderate number of errors. Using a single aggregate value to compute the best route cannot provide good service to both applications simultaneously.
A reactive protocol is a MANET routing algorithm that uses the Hop Count value as the metric for choosing the best route. However, research has shown that the Hop Count value alone is not a good routing metric. In addition, when a MANET routing algorithm is applied on a battery-operated device (e.g., cell phone), it is necessary to consider the effects of multi-hop routing on battery life.
Therefore, there is a need in the art for improved routing protocols for use in mobile ad hoc networks. In particular, there is a need for an improved routing protocol that does not rely on a single metric for determining the best routes from a source MANET node to a destination MANET node.

Apparatus and method for determining aggregated link costs in a mobile ad hoc network


An optical imaging method based on a feedback principle in which the specific scan pattern is adapted according to the shape of the sample. The feedback approach produces nanometer-resolved three dimensional images of very small and moving features in live cells and in a matter of seconds. Images of microvilli in live cultured opossum kidney cells expressing NaPi co-transporter proteins with different GFP constructs and images of cell protrusions in a collagen matrix are produced with a resolution of about 20 nm. Along cell protrusions in three dimensional cellular adhesions could be identified to the extracellular matrix.

1. Field of the Technology
The disclosure relates to the field of optical methods for super-resolution imaging of live cells, specifically to an optical imaging method based on a feedback principle in which the specific scan pattern is adapted according to the shape of the sample which produces nanometer-resolved three dimensional images of very small and moving features in live cells and in a matter of seconds.
2. Description of the Prior Art
Several optical microscopy techniques have been previously developed to "break" the diffraction limit and to produce nanometer-resolved images. These techniques can be broadly classified as "physical" techniques in which ingenious approaches are used to break the diffraction limit. For example the STEO (stimulated emission depletion) technique uses stimulated emission to reduce the effective size of the PSF (point spread function). Other methods have used the determination of the center of mass of the fluorescence emission due to single molecules to obtain images with nanometer resolution of cellular features. The PALM and the STORM techniques and their variants also use this approach
Current methods in laser scanning confocal microscopy are based on moving a laser spot in a predetermined pattern, for example in a raster scan path, to obtain the intensity in each point of a plane of focus. These scanning techniques are inefficient however when the features to be imaged are at the nanoscale and sparse since the scanning path crosses the object to be imaged in only a few points. The efficiency of a predetermined scanning pattern, defined as the ratio of the time the laser beam is on a feature with respect to the total time of scanning, further decreases in three dimensions. Live cell structures like protrusions or microvilli are continually remodeled changing shape and position. Current imaging methods, including the super-resolution techniques known as STED and PALM, are inadequate to detect the dynamic of chemical reactions in these tiny three dimensional structures, which occur in the millisecond to second time scale.

Nanometer-scale optical imaging by the modulation tracking (MT) method

The Regents of the University of California


Mesoporous calcium silicate compositions for controlled release of bioactive agents and methods for producing such compositions are disclosed herein. In one embodiment, mesoporous calcium silicate is synthesized by acid modification of wollastonite particles using hydrochloric acid. A hydrated silica gel layer having abundant Si—OH functional groups can be formed on the surface of wollastonite after acid modification. Bruhauer-Emmett-Teller (BET) surface area increased significantly due to acid modification and, in one arrangement, reached over 350 m2/g. Acid modified mesoporous calcium silicate compositions show a higher ability to adsorb protein compared to unmodified particles and demonstrate controlled release kinetics of these proteins.

Bone forming growth factors, such as bone morphogenetic protein (BMP) and transforming growth factor (TGF-β) have been widely investigated in orthopedics and tissue engineering with great success. These molecules stimulate bone cell growth and differentiation, and accelerate bone tissue regeneration [1, 2]. To prevent possible side-effects in non-target tissues, delivery of growth factors are performed in a controlled manner. For example, growth factors can be sequentially administered in a manner that mimics the time profile of the healing process [2-4].
Biomaterials, such as biopolymers [5-8], inorganic ceramics [9-18] and their compositions [19, 20], can act as reservoirs for deliverable protein and other bioactive molecules if they demonstrate a high capacity for protein adsorption while preserving the protein structure and biological activity over time. In addition to controlled protein delivery, carriers that are biocompatible with bone tissue can also serve as a scaffold for new tissue formation during normal tissue repair or healing. Furthermore, such biomaterials should be non-immunogenic and, in many cases, biodegradable once tissue regeneration is complete [2, 3]. Typically, biopolymers are not bioactive because they do not chemically bond to living bone. Inorganic bone replacement materials, including calcium phosphate ceramics [10, 11], bioglass and bioglass-ceramics [12], silica gel [13,14], and calcium phosphate cements [15-18], have been extensively investigated due to their excellent bioactivity with bone tissue. However, the adsorption capacity of protein on these carriers is limited to outer surface adsorption. Moreover, protein release from these materials displays a burst release (e.g., rapid release) due to the weak bonding between protein and the carriers.

Mesoporous calcium silicate compositions and methods for synthesis of mesoporous calcium silicate for controlled release of bioactive agents


In an integrated circuit comprising discrete circuit modules, dedicated test signals for the individual circuit modules are multiplexed with operational signals to the external pins of the integrated circuit. When not in a test mode, the external pins of the chip are coupled to the normal operational signals. When a test mode is commanded, certain external pins are coupled to test signals that are not otherwise available off-chip in accordance with the contents of a test register.

As the complexity of integrated circuits increases, it becomes increasingly difficult to comprehensively test the performance of such circuits. These circuits are normally subjected to automated test procedures whereby various combinations of input signals are applied to the chip and output signals are monitored for expected results. To exhaustively test a complex integrated circuit, numerous combinations of input signals or test vectors are required.
Integrated circuits that comprise numerous circuit modules, such as may be created using cell based design, will have numerous module pin signals which would make the pin count for the circuit prohibitively large. Thus, it may not be practical to bring all such signals off chip. Furthermore, integrated circuits have numerous signals that are not accessible off-chip and it may not be possible to exercise all internal states of the circuit using input-to-output test techniques of the prior art.
An alternative approach is to mechanically probe the signal traces on the circuit chip. However, such a technique is impractical outside of a laboratory and could not be utilized for production quantities of integrated circuits.
A primary objective of the present invention is to provide a means for gaining access to internal signals of an integrated circuit for test purposes.
It is another objective of the present invention to gain such test access to internal signals without increasing the pin-outs of the integrated circuit chip.

Method and apparatus for mapping test signals of an integrated circuit


An apparatus generates three dimensional image data. Specifically, image data for a first image layer and image data for a second image layer is generated together with depth indication data for at least one of the first image layer and the second image layer. The two image layers may be used e.g. for dependent images, such as foreground and background images, or for independent images, such as initial and final images for a cross fade. A dependency indicator which is indicative of a dependency between the image layers is generated and included in the three dimensional data. A rendering device may receive the three dimensional data and generate images for a three dimensional representation based on the data. The three dimensional processing is adapted in dependence on the dependency indicator. For example, interlayer de-occlusion processing may only be performed if the dependency indicator is indicative of dependent image layers.

Three dimensional displays add a third dimension to the viewing experience by providing a viewer's eyes with different views of the scene being watched. A popular approach for representing three dimensional (3D) images is to use one or more two dimensional (2D) images plus a depth representation providing information of the third dimension. Such an approach provides a number of advantages including allowing 3D views to be generated with relatively low complexity and providing an efficient data representation thereby reducing e.g. storage and communication resource requirements for 3D image (and video) signals. The approach also allows 2D images to be generated with different viewpoints and viewing angles than the 2D images that are included in the 3D image data. This can specifically be used to generate an additional image so that two images corresponding to a viewing angle of a viewers two eyes is provided. This may be used to provide 3D perception, i.e. to provide a 3D image.
A drawback of representing a 3D image by a single 2D image and associated depth information is that it does not include information about background image areas that are occluded by foreground objects. Accordingly, if the scene is rendered for a different viewpoint, no information can be revealed behind the foreground objects. Accordingly, it has been proposed to use multi-layer image and depth representations comprising a plurality of two dimensional images (e.g. a foreground image and a background image) with associated depth information. A description of how to render new views from such information can be found in Steven J. Gortler and Li-wei He, Rendering Layered Depth Images, Microsoft Technical Report MSTR-TR-97-09, and in for example United States Patent Application US20070057944.
In approaches using more than one layer, it has been proposed to allow layers to be semi-transparent. In the field of computer graphics, such an approach is for example described in Norman P. Jouppi and Chun-Fa Chang, "An Economical Hardware Technique for High-Quality Antialiasing and Transparency", Proceedings of Eurographics/Siggraph workshop on graphics hardware 1999. Such an approach allows semi-transparent materials to be visualised (e.g. water, smoke, flames) and also allows improved anti-aliasing of edges of objects at different depths. Specifically, it allows for a more gradual transition of edges. Thus, the transparency may not just be used for representing semi-transparent objects, but may also allow anti-aliasing of the edges of the foreground objects by making an edge semi-transparent such that a transparency value represents how much of a pixel should be foreground and how much of the background should visible. An example of such an approach can be found in: C. Lawrence Zitnick Sing Bing Kang Matthew Uyttendaele Simon Winder Richard Szeliski, "High-quality video view interpolation using a layered representation", in Proceedings of Siggraph 2004.
However, although such approaches may provide improved 3D image perception and images in many scenarios, they tend to have some disadvantages. In particular, the image quality and 3D perception may be suboptimal in some scenarios and the 3D processing may introduce distortions and image degradations, such as visible artefacts.
Hence, an improved approach for 3D image data processing would be advantageous and in particular an approach allowing increased flexibility, improved image quality, facilitated implementation and/or improved performance would be advantageous.


Method for performing timing closure of integrated circuits in the presence of manufacturing and environmental variations. The starting design is analyzed using statistical static timing analysis to determine timing violations. Each timing violation in its statistical canonical form is examined. In a first aspect of the invention, the canonical failing slack is inspected to determine what type of move is most likely to fix the timing violation taking into account all relevant manufacturing and environmental variations. In a second aspect of the invention, pre-characterized moves such as insertion of delay pad cells are evaluated for their ability to fix the timing violation without triggering timing, and the best move or set of moves is selected.

Timing closure refers to the process of making changes to an initial integrated circuit design to get rid of timing violations. Typically, a static timing analysis software tool is employed to determine all the timing violations. Changes are made to improve the design and static timing analysis is iteratively invoked until all violations have been addressed. These changes are also referenced to as "moves." Making moves to improve the timing of the circuit will be referred hereinafter as optimization.
Examples of changes made to achieve timing closure include buffer insertion, gate sizing, assignment of wires to metal layers, pin swapping and pad cell (or "delay cell") insertion. Some moves improve the timing characteristics of the design and others make timing worse, so that the timing closure process is often "timing driven", meaning that the static timer is queried to determine which changes improve the design and which do not. In general, the iterative implementation of changes in a timing-driven manner can be made manually by a designer, or by a software Electronic Design Automation tool.
For an integrated circuit to be free of timing violations, numerous timing tests associated with the circuit design must be met. Associated with each timing test is a timing slack, which is computed by the static timing analysis (STA) tool. Positive slack signifies that the timing is met, and negative slack implies a timing violation. The process of optimization consists of iteratively making moves until a positive slack is achieved for each timing test. STA software tool typically supports slack queries to aid in the optimization process.
Manufacturing and environmental variations significantly complicate the timing closure. Manufacturing variations will be referenced hereinafter as "process variations." Variations are handled by an STA program by utilizing either multi-corner methods or statistical embodiments. The corner is a unique setting of process and environmental variables, and multi-corner timing implies that timing analysis is repeated at multiple corners and the worst (i.e., smallest) slack across all corners must be positive for closure. In a statistical timing, some or all variations are treated statistically, i.e., as probability distributions. The resulting statistical slack is "projected" to obtain the worst slack across the space of manufacturing and environmental variations, and the statistical slack must be positive to achieve closure.
For conventional optimization methods, timing are queried after every move, and the moves that result in an improved timing are accepted, while those that worsen the timing are rejected. One timing test is targeted at a time and various moves are tried in an iterative manner until either a positive slack is achieved or the present timing test is skipped in order to make progress on other timing tests. In this conventional method, the static timing analysis software typically updates timing information incrementally to answer timing queries made after each move. Incremental timing implies that the timing tool performs minimal computation in order to answer the timing query after a circuit change. In another conventional method, a plurality of moves is "batched up" and timing is then checked after applying all the changes to the design.
Both conventional methods suffer from significant drawbacks. In the aforementioned first method, the computer run time of updating timing and answering timing queries after each design change can be prohibitive, because millions of changes can be attempted in a single optimization run and timing queries must be answered every single time. A second drawback is that each proposed change provides a different amount of benefit at different process corners. If the multi-corner or statistical timing must be repeated to evaluate the benefit of each proposed change, the computational burden can become overwhelming. If single corner timing is used to guide the optimization, then the timing violations at other corners may not be addressed and may in fact get worse. Thus, the optimization may improve the timing at the single corner considered, but make timing worse at other corners, and may "ping pong" back and forth without making any real progress. This drawback of prior art embodiments highlights the difficulty of timing closure in the presence of variations. A third weakness of prior art methods is that there is no easy way to decide which type of move is most effective in any given situation.
Programs will typically try either buffering or sizing or layer assignment or other types of moves in different heuristically selected orders to see which yields most benefit. The selected order applies to all timing violations. Presently, violations that benefit from a different order of moves cannot be achieved with conventional methods. As a result, conventional embodiments are inefficient due to applying the same order of types of moves to all the timing violations. This inefficiency manifests itself as longer run times and sub-optimal results in terms of timing, power and area of the integrated circuit. The heuristic ordering of moves makes the optimization results inconsistent, untrustworthy and subject to chance.
Other conventional methodologies also suffer from another significant drawback. As moves are accepted and design changes are implemented, the timing information corresponding to the changed circuit is not available. Therefore, there is no clear timing guidance for the rest of the moves in the "batch" of moves. As in this conventional embodiment, there is no guarantee that various corners or portions of the space of manufacturing variations will have fixed timing violations. As previously stated, there is no intuitive approach for which type of change will be most beneficial in any given situation.
Conventional embodiments often use exhaustive methods. For example, all sizes of a buffer may be inserted one at a time with a timing query following each insertion to determine the best size. In many instances, they are limited by the use of heuristic methods to optimize moves in various different orders to fix the timing violations. For instance, conventional methods may use delay pad cell insertion first, followed by buffering, gate sizing and finally pin swapping. The heuristic order is applied to all the timing violations, which is wasteful because different timing violations benefit from different types of moves, and trying moves in the same order across all the timing violations, which is inefficient.
Thus prior art embodiments have a relatively low fraction of moves that are accepted, and they create sub-optimal designs by inserting buffers, in instances where gate sizing may be a more efficient way to fix a particular timing violation.
In view of these and other reasons, conventional embodiments have failed to achieve timing closure in the presence of manufacturing and environmental variations, and suffer from significant drawbacks.
Accordingly, methods have been proposed in the literature, but generally have not been practiced in industry. Thus, there is a need for methods capable of providing a statistical optimization of integrated circuit designs.

Method for achieving an efficient statistical optimization of integrated circuits

International Business Machines Corporation


Provided is a separator including: a first header into which gas flows from an exterior of the separator; a core connected to the first header and including at least one flow path configured to cool the gas flowing from the first header through the flow path; and a second header connected to the core and configured to guide the gas passing through the at least one flow path to the exterior of the separator, wherein the second header includes at least one first baffle provided on an inner wall of the second header and configured to remove moisture from the gas.

1. Field
Apparatuses consistent with exemplary embodiments relate to an apparatus, and more particularly, to a separator.
2. Description of the Related Art
The durability and operational efficiency of an apparatus using gas as a fuel, such as a multistage compressor, a turbine, or a gas turbine, may be improved by efficiently removing moisture from the gas. The moisture included in the gas may be removed by providing a separator in an area where the gas flows into the apparatus. Also, the separator may be used to condense and extract the moisture from the flowing gas and supply dry gas to the apparatus described above. The gas passing through the separator may have a differential pressure. In addition, a size and volume of the separator may have to be increased to improve the separation efficiency according to the type of the apparatus. Also, a large separator with an increased size and volume may be heavy and require a high manufacturing cost. Therefore, it is important to efficiently manufacture and reduce the size of the separator.


A folding plant and vine supporter that has rings of different diameters which are attached to four identical vertical supports. When fully assembled, the rings are horizontal or situated perpendicular to the vertical supports. The vertical supports provide four equidistant focal points, which maintain the rings in a horizontal position. Each ring rests at a level along the vertical support. The ring with the smallest diameter is the lowermost ring or the ring closest to the ground. The ring with the largest diameter is the uppermost ring or the ring furthest away from the ground. The rings increase in diameter as the distance from the ground increases. Each ring has a ridge around the inside circumference to stabilize the plant support in the assembled position. The rings resist upward movement.

One of the most versatile tools in gardening is plant stands, which are easy to assemble, and tough enough to provide adequate support to a plant throughout the growing season. There are existing plant supports on the market that support vines in a vertical fashion, but as will be shown; none meet all the needs of avid gardeners. Many existing art, are intended only for use with a potted plant, and would not function as a support for earth embedded plants. Existent are many tie off methods, in which the vine is tied to a stake or pole. An obvious problem with this method exists in the possibility of cutting off chloroform flow to the ends of the plant that is tied, thereby, cutting down production of the plant, and possibly killing the plant all together. The tie off methods also do not give the support and protection, as do enclosure supports. Additionally many existing enclosure supports are not collapsible, and are not intended to last more than one growing season.
U.S. Pat. No. 4,026,068 issued to Tepper on May 31, 1977 shows a tomato support stand. Tepper's invention is unlike the present invention because it is not conical; in shape, does not have rings as a support to the vines, and is not collapsible.
U.S. Pat. No. 4,519,162 issued to Stuckey on May 28, 1985 shows a vertical plant support with horizontal hoops. Stuckey's invention is unlike the present invention because it has only one vertical support, the rings are of the same diameter, it is made of steel, and it is not collapsible.
U.S. Pat. No. 4,631,861 issued to Wuthrich on Dec. 30, 1986 shows a combination pot and trellis. Wuthrich's invention is unlike the present invention because it is a planting pot with attached trellis, is not collapsible, and is not functional as a separate trellis for plants embedded in the earth.
U.S. Pat. No. 4,860,489 issued to Bork on Aug. 29, 1989 shows an adjustable plant support. Bork's invention is unlike the present invention because there are two pole supports and a spiral internal member, instead of several separate rings, the spiral member is the same diameter throughout Bork's invention, and it is not conical in shape.
U.S. Pat. No. 4,914,857 issued to Dodgen on Apr. 10, 1990 shows a tomato vine-supporting device. Dodgen's invention is unlike the present invention because it is not conical, is not collapsible, and has an exterior ring that circumscribes the invention while hanging on hook extensions.
U.S. Pat. No. 760,789 issued to Kunzman on May 24, 1904 shows a plant support. Kunzman's invention is unlike the present invention because it is not conical in shape, the rings are not one congruous piece, and the rings are of the same diameter to one another.
U.S. Pat. No. 5,323,559 issued to Allman on Jun. 28, 1994 shows a collapsible plant support. Allman's invention is unlike the present invention because it is not conical in shape, has polygonal rings, instead of circular rings, and the polygonal members are all of the same diameter to one another.
U.S. Pat. No. 5,440,834 issued to Kleinert on Aug. 15, 1995 shows a single fastener plant support ring for use with a T-shaped upright post. Kleinert's invention is unlike the present invention because, it has a single ring attachment, is not collapsible, is not conical in shape, and does not provide a plurality of support points.
U.S. Pat. No. 5,450,691 issued to Christie, et al. on Sep. 19, 1995 shows a support for plants and the like. Christie's invention is unlike the present invention, because it has a ring attachment at the uppermost point that is not congruous with other rings, it is not collapsible, and although it is conical in shape the lesser diameter is intended to be at the top of the invention instead of the base as in the present invention.
U.S. Pat. No. 5,544,446 issued to Benson on Aug. 13, 1996 shows a collapsible and adjustable plant support and protector. Benson's invention is unlike the present invention because it is a polygonal enclosure for plants, does not have a conical shape, does not have rings increasing in diameter and is not intended as a vine support, but rather a pen enclosure for a variety of plants.
U.S. Pat. No. 5,595,019 issued to Foreman on Jan. 21, 1997 shows a plant support. Foreman's invention is unlike the present invention because it has a center support that protrudes from the ground, there is a larger gap in the center region between the two middle rings of the support, and it is not collapsible in nature.
Therefore a need has been established for a conical support for vines, and the like, that is collapsible, and embeds in the earth.


A door opening and closing apparatus for a vehicle includes a support member interposed between a body and a door of the vehicle and expandable and shrinkable so as to enable holding of the door at an open position. The support member includes a first housing unit having an electric motor driving mechanism and a second housing unit integrally joined to the first housing unit. The second housing unit includes an outer cylinder, an inner cylinder disposed concentrically in the outer cylinder and engaged with the outer cylinder by fitting engagement in a movable manner relative to the outer cylinder, a spindle connected to the electric motor driving mechanism, and a spindle driving mechanism for moving the inner cylinder relative to the outer cylinder in an axial direction due to rotation of the spindle.

Technical Field
The present invention relates to a door opening and closing apparatus for a vehicle.
Related Art
A known door opening and closing apparatus for a vehicle includes a support member interposed between a back door which opens or closes a rear opening portion of the vehicle and a body of the vehicle, and can hold the back door at an open position by the support member (see JP 2014-101637 A, for example). The support member includes an inner cylinder and an outer cylinder which are disposed concentrically and are slidably engaged with each other by fitting engagement. In the inside of the inner cylinder, elements such as an electric motor which generates a drive force, a speed reduction gear mechanism which decelerates the rotation of the electric motor are housed. In the inside of the outer cylinder, elements such as a spindle, a spindle nut, a guide tube, a coil spring, a push rod are housed. The spindle is rotatably driven by the electric motor. The spindle nut is threadedly engaged with the spindle. Due to the rotation of the spindle, the inner cylinder is moved in an axial direction together with the spindle nut and the push rod. The coil spring biases the push rod and the outer cylinder in an axial direction. The guide tube houses the coil spring therein and slidably guides the push rod.
In the door opening and closing apparatus for a vehicle which includes such a support member, in a manipulation of opening the back door in a closed state, the electric motor is driven in response to such an open manipulation. The rotation of the electric motor is transmitted to the spindle while the rotation of the electric motor is decelerated by the speed reduction gear mechanism, and the spindle is rotatably driven at a decelerated speed. When the spindle is rotated, the spindle nut is moved (pushed out) in an axial direction together with the push rod and the inner cylinder. As a result, the inner cylinder is moved relative to the outer cylinder so that the support member is extended whereby the back door is moved toward an open position. At this stage of the operation, the back door can be stopped at any desired position due to a load balance among a biasing force of the coil spring, a biasing force of a gas spring, slide resistances of the speed reduction gear mechanism and the electric motor and the like. Accordingly, the back door is held at the open position. When a manipulation of closing the back door at the open position is performed, the electric motor is rotated reversely in response to such a closing manipulation and the spindle is also rotated reversely and hence, the spindle nut is moved in a reverse direction together with the guide tube and the inner cylinder. As a result, the inner cylinder is moved (pulled) relative to the outer cylinder against a biasing force of the coil spring whereby the support member is shrunken and the door is driven to the closed position.
Conventionally, with respect to this type of support member, the electric motor which is required to possess a watertight structure for preventing leakage of electricity and the spindle, the coil spring and the like which are not required to possess high watertight property are incorporated together in the support member. When the support member adopts such a structure, it is necessary to assemble all parts through consistent manufacturing steps. Accordingly, conventionally, a manufacturing step control of this type of support member is not easy so that productivity of the support member is low.
Further, an opening/closing angle, a mounting point and the like differ for respective vehicle types and hence, when it is necessary to provide support members having different lengths for the respective vehicle types, it is necessary to change the structure of the whole support member. Such a change of the structure pushes up a cost of the support member.

Door opening and closing apparatus for vehicle


Digital electronic apparatus for performing chain-serial matrix multiplications using a single pipeline multiplier supplies elements of the multiplicand and multiplier matrices to the digital memory from first and second memories. Each product matrix is temporarily stored in a third memory until such time as it is used to write the first memory for the next matrix multiplication in the series. This procedure avoids overwriting the first memory when its data are still required for application to the pipeline multiplier.

The synopsis "Matrix Algebra" appearing on pages 1090-1097 of Chapter 37 of the fourth edition of REFERENCE DATA FOR RADIO ENGINEERS published by International Telephone and Telegraph Corporation, 67 Broad Street, New York 4, N.Y. is hereby included herein by reference. Chain-serial operation, insofar as matrix multiplication is concerned, can be sustained in two forms. In a first form of chain-serial matrix multiplication, an m row by n column original matrix is continuously updated by using it as the left hand multiplicand term in each of a succession of matrix multiplications with a right hand multiplier term that is the next one of a succession of n row by n column matrices. In a second form of chain-serial matrix multiplication, an n row by m column original matrix is continuously updated by using it as the right hand multiplier term in each of a succession of matrix multiplications with a left hand multiplicand term that is the next one of a succession of n row by n column matrices. Each successive matrix multiplication in the first form of chain-serial operation can be described by the matrix equation (a.sub.ij).multidot.(b.sub.ij)=(c.sub.ij) written more explicitly as follows: ##EQU1## In the matrix equation just above ##EQU2## where i and j are the numbers of the row and of the column respectively of each element in the (a.sub.ij), (b.sub.ij) and (c.sub.ij) matrices. Note that (b.sub.ij) is a square matrix, having n rows and n columns. The matrices (a.sub.ij) and (c.sub.ij) have m rows and n columns; they may be square with m equal to n, or they may not be.
Each successive matrix multiplication in the second form of chain serial operation can be described by the matrix equation (d.sub.ij).multidot.(e.sub.ij)=(f.sub.ij) written more explicitly as follows: ##EQU3## In the matrix equation just above ##EQU4## where i and j are the numbers of the row and of the column respectively of each element in the (d.sub.ij), (e.sub.ij) and (f.sub.ij) matrices. Note that (d.sub.ij) is a square matrix, having n rows and n columns. The matrices (e.sub.ij) and (f.sub.ij) may be square or may not be.
It is known to those familiar with matrix algebra that the product of the transposes of two square matrices transposes to the product of the two matrices multiplied in reverse order from their transposes. That is, (a.sub.ij).multidot.(b.sub.ij) is the transpose of (b.sub.ji).multidot.(a.sub.ji), where (b.sub.ji) is the transpose of (b.sub.ij) and where (a.sub.ji) is the transpose of (a.sub.ij). Accordingly, each step of matrix multiplication to be done in one of the first and second forms of chain-serial multiplication can alternatively be done in the other of those forms by transposing the matrices entering the reverse-order matrix multiplication and transposing the resulting matrix product. One can identify (a.sub.ji) with (e.sub.ij) and can identify (b.sub.ji) with (d.sub.ij). Then (c.sub.ij)) equals (f.sub.ji), the transpose of (f.sub.ij). While the pipeline circuitry in a chain-serial matrix multiplier may be designed particularly to suit one of the two forms of chain-serial matrix multiplication that is self-sustaining, it can be operated so as to provide the other form of chain-serial multiplication, by transposing rows and columns of the matrices involved in each matrix multiplication and by performing that multiplication in reverse order, then transposing the result.
In chain-serial matrix multipliers that embody the invention in various of its aspects, the multiplications of pairs of elements involved in the matrix multiplication are (except possibly for trivial multiplications) performed by a single multiplier. First and second streams of matrix elements are supplied to the multiplier and an ensuing accumulator, which together perform pipelined calculation of elements of the product matrix. The first stream of matrix elements is generated by scanning a first matrix as stored temporarily in and read subsequently from a first memory. This scanning may be done in row major order (i.e., by rows of elements) or in column major order (i.e., by columns of elements). The second stream of matrix elements is generated by scanning a second matrix as stored temporarily in and read subsequently from a second memory in the opposite order that said first matrix is scanned.
Where desired, transposition of the first and second matrices is easily arranged for, by transposing the order the first and second memories are written. Owing to the fact that matrices may be so easily transposed in the operation of the chain-serial matrix multipliers of the invention, a general term referring to a row or a column of a matrix is convenient to have in claiming these chain-serial matrix multipliers. Accordingly, a "line" of a matrix refers to a row of a matrix the element samples of which flow in row major order through a pipeline; and a "line" of a matrix also refers to a column of a matrix the element samples of which flow in column major order through a pipeline.
A type of memory that is convenient to use in pipelined operations involving arrays of data, such as those involved in finite-impulse-response filtering in one or more dimensions, is so-called "serial memory". Serial memory comprises a plurality of clocked delay elements, connected one after the other in cascade connection to form a tapped delay line. Each delay element exhibits a one clock interval delay between its input and output ports and provides for temporary storage of a respective sample. A multiplexer selectively connects to the input port of the first delay element in the cascade connection (i.e., the initial tap of the delay line) either a sample being written into the serial memory or a sample supplied from the output port of the last delay element in the cascade connection. The latter selection by the multiplexer provides for a circulating delay line, or ring connection of the delay elements, which sustains continuing storage of the samples written into the serial memory. Samples may be read from the serial memory at any tap along the tapped delay line (e.g., the initial tap or the final tap).


Examples of a carry chain for performing an operation on operands each including elements of a selectable size is provided. Advantageously, the carry chain adapts to elements of different sizes. The carry chain determines a mask based on a selected size of an element. The carry chain selects, based on the mask, whether to carry a partial result of an operation performed on corresponding first portions of a first operand and a second operand into a next operation. The next operation is performed on corresponding second portions of the first operand and the second operand, and, based on the selection, the partial result of the operation. The carry chain stores, in a memory, a result formed from outputs of the operation and the next operation.

The term single instruction, multiple data (SIMD) describes computers with multiple processing elements that perform the same operation on multiple data elements simultaneously. Most modern CPU designs, such as ARMv8, include SIMD instructions in order to improve the performance of multimedia use, like graphics or audio processing; SIMD instructions can also be used to accelerate digital signal processing in wireless communications systems and cryptographic applications. The data element width or size can vary, depending on the application; 8-bit, 16-bit, 32-bit, and 64-bit elements are commonly supported.


An electronically controlled braking system includes at least one central control unit, the at least one central control unit operable to assign identifiers during an identifier assignment routine, a first distributed electronic control unit, a second distributed electronic control unit and a control network over which the at least one control unit, the first distributed electronic control unit and the second distributed electronic control unit are communicable. The identifiers include a first identifier and a second identifier, one of which is assigned to the first distributed electronic control unit and the other which is assigned to the second distributed electronic control unit at least in part based upon the timing of identification signals generated by the first distributed electronic control unit and the second distributed electronic control unit reaching the central control unit via the control network.

Traditional braking systems for motor vehicles include conventional hydraulic or pneumatic brakes associated with two or more wheels of the vehicle. Such conventional brakes are actuated by pressurized fluid or compressed air. When actuated, the brakes exert a force on a disk or drum which spins in conjunction with the wheel of the vehicle in order to create frictional forces which resist rotation of the wheel. Traditionally, control signals have been transmitted to each of the brake system's actuators mechanically, or by a hydraulic or pneumatic control circuit. However, it has more recently been proposed to employ a centralized control unit to generate electronic control signals and to use such electronic control signals to control actuation of a vehicle's brakes. This type of electronic control scheme has become even more prevalent in view of modern brake systems which now often include not only conventional hydraulic or pneumatic brake actuator functionality, but also supplemental electronic functions such as antilock protection (ABS) and/or electronic braking force distribution (EBV) between the front and rear axles.
In some brake systems employing electronic control, it is known to employ distributed electronic control units to receive electronic control signals from the central control unit and to control the actuation of brake components in response thereto. Such distributed electronic control units may be associated with each individual brake component and/or may be associated with two or more brake components (for example, one distributed electronic control unit may be associated with both brake actuators on a single axle). In either case, the central control unit generates and transmits electronic control signals to each of the distributed electronic control units via one or more control networks, and the distributed electronic control units, in response to such electronic control signals, cause the brake component(s) with which they are associated to actuate. Such a system, however, suffers from a number of disadvantages. One disadvantage of such a system is that all brake components must be responsive to a single control signal, as there is no way for control signals intended for controlling only one brake component or group of brake components to be differentiated from those control signals intended for controlling another brake component or group of brake components.
This disadvantage is obviated by certain brake systems which employ a scheme for the recognition of the physical position of each individual brake component or each group of brake components. Known schemes of this type generally involve the manual programming or setting of a unique identifier in the memory of each distributed electronic control unit, which identifier is indicative of the physical position of the brake component or group of brake components with which the distributed electronic control unit is associated. In these systems, each of the control signals generated by the central control unit includes an indication of the identifier of the distributed control unit for which that control signal is intended. When each distributed electronic control unit receives the control signals, it determines, based upon the identifier contained in each control signal, whether that particular control signal is intended for it, and if so, responds accordingly. If that particular control signal is not intended for it, the control signal is ignored.
For example, the distributed electronic control unit associated with a brake component on the right front wheel of a vehicle may be programmed with an identifier indicative of such a position. Thus, when control signals are received from the central control unit, the distributed electronic control unit associated with a brake component on the right front wheel of the vehicle responds to the control signal only if that control signal includes the appropriate identifier (if not, the control signal is ignored). Distributed electronic control units associated with groups of brake components (e.g., the brake components on a vehicle axle) may be similarly programmed and controlled.
While the above-described control scheme allows each individual brake component or each group of brake components to be independently controlled by the central control unit, it does still suffer from a number of disadvantages. More specifically, the manual programming or setting of the unique identifier in the memory of each distributed electronic control unit may create compatibility problems, which may lead the brake system to operate inefficiently or even fail completely. This is true because each time a distributed electronic control unit is replaced and/or its location in the system is changed, the distributed electronic control unit must be manually programmed or reprogrammed. For example, if a brake component including a distributed electronic control unit failed and was replaced, the new brake component must be programmed with the identifier appropriate for its location. If the service technician overlooks such programming or programs the distributed electronic control unit incorrectly, that brake component may not respond at all, or may respond at inappropriate times.
A similar problem may occur if the locations of two or more brake components are switched. For example, if the left front and right front brake components are switched during servicing and not reprogrammed, the control signals generated by the central control unit intended to control the left front brake component would actually be controlling the right front brake component, while the control signals intended to control the right front brake component would actually be controlling the left front brake component. This could cause a hazardous condition during vehicle operation. Moreover, even if programming and/or reprogramming is carried out correctly, the process may be time-consuming and would likely require additional knowledge on the part of service technicians who service the brake system.
What is desired, therefore, is an electronically controlled braking system intended for use with wheeled vehicles which incorporates distributed electronic control units, which allows for the control of each individual brake component or each group of brake components associated with each distributed electronic control unit independently of those associated with other distributed electronic control units, which allows for control signals intended for controlling only one brake component or group of brake components to be differentiated from those control signals intended for controlling another brake component or group of brake components, which does not require that the distributed electronic control unit be manually programmed or reprogrammed each time a distributed electronic control unit is replaced and/or its location in the system is changed, and which employs an identification scheme, the implementation of which is not time-consuming and/or does not require additional knowledge on the part of service technicians who service the brake system.

This patent is subject to a terminal disclaimer.

Brake system with distributed electronic control units


A novel washing method of a washing machine that is capable of performing a washing operation of laundry with reduced washing water consumption and low power consumption and, in addition, sterilizing the laundry is disclosed. The washing method includes supplying washing water containing detergent into a tub to perform a high-concentration washing process, and, before or after the high-concentration washing process, supplying steam into a drum.

Generally, mashing machines are classified into a pulsator type washing machine, a drum of which is mounted in a vertical direction, and a drum type washing machine, the drum of which is mounted in a horizontal direction.
In the drum type washing machine, as the drum is mounted in the horizontal direction as described above, laundry received in the drum is washed by a lifting and dropping operation.
FIGS. 1 and 2 schematically illustrate the structure of a conventional drum type washing machine.
As shown in the drawings, the drum type washing machine includes a machine case 10, a tub 20 mounted in the machine case 10, a drum rotatably mounted in the tub 20, and a driving unit for driving the drum 30.
At the front part of the machine case 10 is formed a laundry inlet hole 11, through which laundry is put into the drum. A door 40 is mounted to the machine case adjacent to the laundry inlet hole 11 for opening and closing the laundry inlet hole 11.
At the inner circumference of the laundry inlet hole 11 is mounted a rim part 50, which accomplishes sealing between the door 40 and the laundry inlet hole 11.
At the bottoms of opposite sides of the outer circumference of the tub 20 are mounted dampers 21, which support the tub 20 in the machine case 10.
A washing water heater 60 for heating washing water is mounted in the lower part of the tub 20.
The driving unit includes a driving motor 71 for driving the drum 30, a belt 72 for transmitting the driving force of the driving motor to the drum 30.


Apparatus, methods and systems to produce a protection voltage are disclosed. The apparatus includes circuitry to deliver a first supply voltage to a plurality of circuits, where the first supply voltage has a first magnitude, circuitry to deliver a second supply voltage to a part of the plurality of circuits, where the second supply voltage has a second magnitude, and circuitry to deliver a protection voltage to the part of the plurality of circuits when the second supply voltage is LOW and the first supply voltage is HIGH. The protection voltage has a magnitude that is a fraction of the magnitude of the first supply voltage. The apparatus includes circuitry that causes the delivery of the second supply voltage to the part of the plurality of circuits when the second supply voltage is turned HIGH subsequent to the second supply voltage being LOW when the first supply voltage is HIGH.

An input/output (I/O) interface is a system that enables the transfer of information between internal storage and external I/O devices. In particular, the I/O interface acts as an interface between the central processing unit (CPU) and memory and to one or more peripherals. In many high-speed I/O circuit designs, several different power supply voltages may be used in order to meet a particular I/O specification. For example, some conventional designs use supply VDDQ ranging from 1.2 volts (v) to 1.5v and core logic positive supply (VDD) ranging from 0.8v to 0.9v. Designing I/O circuits that use different supply voltages can present challenges. Such designs can include the use of both high-voltage transistors and low-voltage transistors.
The figures are not to scale. Wherever possible, the same reference numbers will be used throughout the drawing(s) and accompanying written description to refer to the same or like parts.

Methods and apparatus to generate a circuit protection voltage

Hewlett Packard Enterprise Development LP


Systems, methods, and software applications for enabling the collaboration on a project plan are described in the present disclosure. A computer readable medium is configured to store instructions that are executable by a processing device. According to one embodiment, among many, the computer readable medium includes logic adapted to enable a member of a project team to submit a proposal for modifying a current project plan to a project manager. The computer readable medium also includes logic adapted to enable the project manager to accept or reject the proposal for modifying the current project plan. Various team members make changes to a single shared copy of the project plan. The changes can be to a respective team member's section of the plan

In a business environment, a team of people may be grouped together in order to complete a specific project. In order to define how a project is to be completed and how certain goals are to be accomplished, a project plan can be utilized. The project plan lists different tasks and assigns those tasks to different people on the team. Usually, a project plan can contain tasks that are related hierarchically, where some tasks must be completed before other tasks can be started. Also, a project plan can often have dates or times when certain tasks are to be completed in order for the project to run smoothly for the entire team.

Enabling collaboration on a project plan


A method and apparatus for efficiently maintaining XML index structures in a database system is provided. Instead of immediately synchronizing the XML index structures for every change submitted to the database, the changes are stored in a PENDING table. The XML index structures are periodically synchronized, based on certain criteria or at a user's discretion. Between synchronizations, the XML index structures may be in stale mode or current mode. If in stale mode, a request to access indexed information will use the XML index structures to find the desired indexed information without checking the PENDING table. If in current mode, a request to access indexed information will search the XML index structures and also check the PENDING table to determine 1) whether the indexed information has been updated or deleted, and 2) whether information relevant to the request has been inserted or updated.

XML (i.e. eXtensible Markup Language) is becoming the standard format for representing semi-structured data which are then stored and managed within database systems. All the major database systems are being extended to natively support XML data. Standard query languages such as XPath and XQuery are typically used to query these document collections. Specialized index structures, such as an XML index, have been developed to improve the performance of query operations. An XML index stores entries corresponding to all the nodes of a set of indexed XML documents. Each entry tracks information about the corresponding node, such as the path to the node, the value of the node (if the node is a leaf node) and some information about the hierarchical position of the node (such as a dewey style order key). Additional secondary indexes may also be created to improve lookup of entries based on path, value, etc.
However, creating an XML index adds significant overhead to the cost of inserting and updating documents within the system. When a new document is inserted, or an existing document is updated, entries corresponding to all the affected nodes of the document need to be generated and inserted into the index structure. In the case of an update operation, all the entries corresponding to the nodes of the old document need to be deleted and entries corresponding to the nodes of the updated document need to be added. Further, any secondary indexes need to be suitably updated.
To ensure that the index contains accurate information, changes made to the index and any secondary indexes are performed in the same atomic transaction as the change to the underlying data. In other words, the index and the underlying data are maintained synchronously. These index maintenance operations could result in a significant slowdown of the overall insert/update operation. This problem is especially critical for high-throughput and low latency requirements of On-Line Transaction Processing (OLTP) applications.
This presents a significant dilemma. On the one hand, although significant overhead is incurred, the index must be updated in order to be useful. On the other hand, simply eliminating the index is not suitable because doing so would severely degrade query performance.
Based on the foregoing, there is a clear need to provide a mechanism for retaining the benefits of an XML index without incurring significant overhead during inserts and updates of documents.
The approaches described in this section are approaches that could be pursued, but not necessarily approaches that have been previously conceived or pursued. Therefore, unless otherwise indicated, it should not be assumed that any of the approaches described in this section qualify as prior art merely by virtue of their inclusion in this section.

Mechanism for efficient maintenance of XML index structures in a database system


Methods, systems, and apparatus, including computer programs encoded on a computer storage medium, for annotating search results. In one aspect, a method includes receiving a notification request, the notification request for identifying when to send a notification from the server system to a device, the notification request includes a device identifier, message text, and a location identifier. The method also includes receiving information corresponding to a current location of the device and information corresponding to an identity of the device. The method also includes determining that the information corresponding to the current location corresponds to the location identifier of the device from the notification request and that the information corresponding to the identity of the device corresponds to the device identifier from the notification request. The method also includes sending the notification, including the message text, to the device.

Individuals carry reminders and lists to remind them of tasks that need to be completed. These reminders can be location dependent. For example, a user may carry a grocery list to a grocery store. Parents may bring a wish list to a toy store for holiday and birthday shopping. Individuals use a variety of reminder methods to remind them to perform specific tasks. For example, an individual may tie a string around his finger to help him remember to pick up the dry cleaning or to purchase an anniversary gift.
Cellular telephones and other mobile devices have the ability to send and receive messages over a radio network. These messages can include Short Messaging Service (SMS) messages, e-mail messages, and push notifications.


A system for determining an actual measurement of a structure on a sample using a measurement tool with a calibration standard having measurement sites. A previously measured site on the calibration standard with a first known metric is measured with the measurement tool to produce a first measurement. A newly measured site on the calibration standard, also with a second known metric is measured with the measurement tool to produce a second measurement. A calibration factor for the measurement tool is computed by comparing the first measurement to the first known metric and the second measurement to the second known metric. The structure on the sample is then measured using the measurement tool to produce a precursor measurement. This precursor measurement is adjusted with the calibration factor to produce an intermediate measurement. Then the intermediate measurement is adjusted with the sample composition data to produce the actual measurement.

This invention relates to the field of instrumentation. More particularly, this invention relates to improved data output from an electron microscope.
In the integrated circuit industry, electron microscopes are central to microstructural analysis of integrated circuit components. The quality of the finished integrated circuit is highly dependent on the measurement and control of an integrated circuit""s critical dimensions. Thus, it is very important to ensure that the critical dimension measurements received from the electron microscope are precise and accurate.
Typically, in critical dimension analysis of an integrated circuit component the electron microscope measures the apparent width of a structure when determining its dimensions. The apparent width of the structure is compared to critical dimension specifications in order to determine the compliance of the integrated circuit component.
Unfortunately, there seem to be disadvantages to using the typical apparatus and method, as the apparent width of a structure as reported by the measurement tool is often different from the actual width of the structure. In addition, the discrepancy between the actual width and the apparent width of the structure seems to fluctuate from sample to sample, and even from day to day. Thus, the integrity of the data derived from such measurements is often called into question, and is difficult to rely on.
In an effort to overcome this problem, some have used a calibration piece having a structure with a known size. The calibration piece is loaded into the measurement tool and measured at regular intervals, such as once each day. The difference between the apparent width and the actual width of the structure on the calibration piece is used as a correction factor for other measurements. Unfortunately, even this procedure tends to not have the desired accuracy in all situations.
Similarly, calibration pieces have been used that are optimized for viewing on an electron microscope, such as tin-on-gold resolution standards. These are used to verify the proper functioning of the electron microscope, and to measure the inherent resolution of the electron microscope. Unfortunately, because the interaction between the electron beam and the calibration piece is very different on such standards in comparison with the interaction between the electron beam and the semiconductor samples to be measured, the data produced is unfortunately of limited use in calibrating the scanning electron microscope for use as a measurement tool.
What is needed, therefore, is a system to improve the precision and accuracy of measurement data obtained from an electron microscope during critical dimension review of an integrated circuit component.
The above and other needs are met by a system for determining an actual measurement of physical properties of a structure on a sample using a measurement tool with a calibration standard having measurement sites. A first site on the calibration standard with a first known metric is measured with the measurement tool to produce a first measurement. A calibration factor for the measurement tool is computed by comparing the first measurement to the first known metric. The structure on the sample is then measured using the measurement tool to produce a precursor measurement. This precursor measurement is adjusted with the calibration factor to produce an intermediate measurement. Then the intermediate measurement is adjusted with the sample composition data to produce the actual measurement.
Thus, rather than naively processing the scan data from the measurement tool to produce a measurement result, a model is preferably applied of both (1) how the electron optics perform, including their deviations from ideality, and (2) how the incident electrons interact with the structure on the sample to produce secondary and backscattered electrons. The properties of the electron optical system are preferably derived from both an analytical model of the optical system and from the measurement data taken on the calibration standard. The actual physical properties of the sample to be measured are then preferably determined using an analytical model of the interaction of the incident beam with the sample and the properties of the electron optical system as determined above. The physical properties are preferably determined by iteratively changing some of the sample parameters that are input to the sample interaction model, such as, but not limited to, feature width, height, sidewall angle, and degree of crystallinity, while other sample parameters are preferably kept constant, such as, but not limited to sample material, etc., to produce an analytically derived profile that most faithfully reproduces the empirically observed profile from the sample measurement. In this sense, the present system preferably goes beyond merely determining a correction factor to be applied to the measurement.
In this manner the system described herein provides improved measurement data by correcting the apparent width of a structure by both a calibration factor, which accounts for any drift in the properties of the measurement tool, and by sample structural and composition data, which accounts for measurement differences due to different materials and structures being measured. By calibrating the measurement tool in this manner, the precision and accuracy of the measurement tool is improved. The calibration standard is preferably of a similar composition of the structure to be measured, allowing for correct measurements of structures consisting of a wide range of materials.
Most preferably, the first site is a previously measured site. Preferably, a newly measured site on the calibration standard with a second known metric is also measured with the measurement tool to produce a second measurement. In this embodiment, the second measurement is also compared to the second known metric, and this information is also used in the computation of the calibration factor.
In various preferred embodiments, the measurement tool is preferably an electron microscope, and most preferably a scanning electron microscope.
Thus, rather than merely producing and using a single calibration factor, one or more properties of the electron optical system are preferably determined, which may include incident spot size, collection angles for secondary and backscattered electrons, depth of focus, etc. These properties are preferably used as the inputs to the optical system model used in the sample measurement. Determining these properties from the calibration sample measurement preferably involves the use of known properties of the sample, such as material or nominal feature dimensions, and the use of the sample interaction model described above.
One purpose of using multiple calibration sites is to preferably distinguish between variations in the calibration scan data, which variation can arise from three different sources, being: 1) variation in the measurement tool over time, 2) variation in the properties of one measurement site to another, and 3) variation due to modification of a single measurement site due to interaction with the measurement tool. The source of the variation is preferably determined by using the measurement results from a single calibration standard, as measured at multiple sites on the calibration standard, and the calibration data gathered from previous readings of calibration standards.
This information can be used to simultaneously determine, for example, that the width of a feature is systematically increased by one nanometer each time it is measured, that the size of the irradiated spot has decreased by a tenth of a nanometer since the last calibration, and that the two added measurement sites on the calibration standard deviate from the nominal size by 0.2 and xe2x88x920.1 nanometers respectively. The determination of each of these example determinations is subject to some statistical uncertainty, which can be determined from the data.
The calibration factor is preferably computed by averaging the difference of the first measurement and the first known metric and the difference of the second measurement and the second known metric. The calibration standard is preferably replaced after each measured site has been measured a given number of times, the given number of times being the number of measurements before the measurement tool causes significant structural changes to the measured sites. The calibration standard preferably has over one hundred measurement sites. The composition data for the measured sites on the calibration standard is preferably substantially similar to the sample composition data. The composition data includes material information and structural information. The structural information is preferably automatically determined by imaging the structure. The material information is preferably automatically determined using material analysis techniques, and most preferably using dispersive x-ray analysis.

Critical dimension scanning electron microscope


A storage means and transportation method for at least one (1) motorcycle is herein disclosed, comprising a platform plate supported at corner locations by rotating casters. A vertical edging is provided along the long edges of the platform to prevent the motorcycle from rolling off. At least one (1) handles is provided along the long edge to aid in manipulating each motorcycle along a floor surface for shipping or storage. At each end portion of the platform are fold-down ramps to allow each motorcycle to be rolled on and off of the platform. Each ramp comprises a support to anchor the platform against the floor surface by lifting the casters slightly off of the floor. Once the platform is secured in this fashion, each motorcycle can be rolled on and off the platform safely. Each motorcycle is covered by a fiberglass hard cover which is attached thereto the platform for long-term storage and/or shipment.

Motorcycles represent an increasing sector of the motor vehicle industry. Their cheap maintenance and fuel costs and avid customer loyalty make such vehicles an ever popular presence on the market. However, the use of such vehicles is not suited for all situations. In some cases motorcycles must be transported from one (1) location to the next without the actual running of the motorcycle. Such transactions often occur between the point of manufacture and the point of sale, during resale of the vehicle when a buyer must travel in order to purchase the motorcycle, and the like.
Various attempts have been made to provide apparatuses which aid in the lifting or transport of motorcycles and the like. Examples of these attempts can be seen by reference to several U.S. patents. U.S. Pat. No. 3,785,517, issued in the name of Brajkovich, describes a motorcycle carrier and chock adapted for securing and transporting a motorcycle via a truck, trailer or similar existing vehicle.
U.S. Pat. No. 5,609,461, issued in the name of Lichtenberg, describes a motorcycle dolly adapted for the local transport of a motorcycle via a rolling apparatus, thus allowing a user to move the motorcycle easily and without need to run the motorcycle.
U.S. Pat. No. 6,575,310, issued in the name of Chamoun, describes a motorcycle lift. The Chamoun apparatus comprises a platform and a mechanism adapted for easily and securely lifting a motorcycle positioned on the platform.
Additionally, ornamental designs for a motorcycle lift exist, particularly U.S. Pat. Nos. D 378,155 and D 495,107. However, none of these designs are similar to the present invention.
While these apparatuses fulfill their respective, particular objectives, each of these references suffer from one (1) or more of the aforementioned disadvantages. Many such apparatuses do not provide protection or storage capabilities. Also, many such apparatuses are suitable only to particular situations. In addition, many such apparatuses are bulky or complicated, making them difficult for a single user to operate. Accordingly, there exists a need for an apparatus adapted for the transport and manipulation of motorcycles without the disadvantages as described above. The development of the present invention substantially departs from the conventional solutions and in doing so fulfills this need.

Storage means and shipping method for motorcycles


A dice-less Craps game is disclosed. The game may be played by dealing, by a dealer, a pair of cards having numerical indicia thereon; determining whether the pair of cards are of a like suit; if the cards are not of a like suit, then adding the numerical indicia together to determine a roll in accordance with the rules of Craps; and if the pair of cards are of a like suit, then arriving at a predetermined result other than a roll in accordance with the rules of Craps. In a further disclosed aspect, the first predetermined result comprises a NoCall. In a further disclosed aspect an apparatus for playing a dice-less Craps game is disclosed, comprising a deck of cards comprising 48 cards having numerical indicia thereon including four like-suited pairs each of Aces, Twos, Threes, Fours, Fives, and Sixes; and a Craps playing surface having additional indicia thereon for wagering on a predetermined result if a pair of dealt cards are of a like suit.

I. Field
The present disclosure relates to card games, and in particular, to a method and apparatus for playing a dice-less Craps game.
II. Background
Craps is one of the most popular dice games played in the casinos of Nevada and New Jersey. Players and spectators alike enjoy the fast-paced action provided by Craps.
However, in some areas of the country dice games are prohibited, while other types of games, such as card games, slot machines, or keno, may be permitted. In such jurisdictions, the absence of Craps on the casino floor may be missed.
Card games designed to replicate dice games are known in the art. For example, the Official World Encyclopedia of Sport and Games, 1979, discloses a craps game that is played with a deck of cards.
The game uses a special deck of 48 cards, consisting of two sets of each of the denominations Ace, 2, 3, 4, 5 and 6. The numerical value of the cards correspond to the indicia on the faces of two dice.
The shooter deals two cards from the top of the deck face up onto the table, constituting a throw. The value of the two cards added together is a xe2x80x9crollxe2x80x9d in the same way as the two dice in dice craps. Play follows the basic rules for Craps, where the first two cards on a come-out give either a natural (7 or 11), a craps (2, 3, or 12), or a point (4, 5, 6, 8, 9 or 10).
After every come-out throw the two cards are shuffled back into the deck by the shooter, and the deck is cut. This happens even if no decision has been reached (e.g. if a point has been thrown), If the shooter then still has to make a point, he deals further throws, but does not shuffle these back into the deck.
This continues until he makes the point or sevens out. The entire deck is then shuffled together and cut. If, on a point, the entire deck is exhausted, the deck is reshuffled and cut by the shooter; he then continues, trying for the same point.
The prior art card games suffers from certain disadvantages, however. For example, since cards are not returned to the deck while the shooter is trying to make a point, the odds of the game are not equivalent to a true Craps dice game since the rolls represented by the used cards are not available.
A novel and improved dice-less Craps game is disclosed. In one disclosed aspect, a method for playing a dice-less Craps game with a deck of cards is disclosed, comprising: dealing, by a dealer, a pair of cards having numerical indicia thereon; determining whether the pair of cards are of a like suit; if the cards are not of a like suit, then adding the numerical indicia together to determine a roll in accordance with the rules of Craps; and if the pair of cards are of a like suit, then arriving at a predetermined result other than a roll in accordance with the rules of Craps. In a further disclosed aspect, the first predetermined result comprises a NoCall.
An apparatus for playing a dice-less Craps game is disclosed comprising: a Craps playing surface for receiving a pair of dealt playing cards, the playing surface having indicia thereon for wagering on a predetermined result other than a roll in accordance with the rules of Craps.
In a further disclosed aspect an apparatus for playing a dice-less Craps game is disclosed, comprising a deck of cards comprising 48 cards having numerical indicia thereon including four like-suited pairs each of Aces, Twos, Threes, Fours, Fives, and Sixes; and a Craps playing surface having additional indicia thereon for wagering on a predetermined result if a pair of dealt cards are of a like suit.
In additional aspect of a disclosed dice-less Craps game, a pair of dealt cards may be drawn from a deck consisting of 48 cards including four like-suited pairs each of Aces, Twos, Threes, Fours, Fives, and Sixes. The dealt pair of cards may be returned to the deck prior to the dealer dealing a subsequent pair of cards.
It is contemplated that the disclosed game may be embodied in computerized gaming equipment.

Method and apparatus for playing a diceless craps game


A cradle assembly for housing and supporting a motor and a gearbox of an electric axle assembly. The cradle assembly has generally an “open air” design which provides a space for a motor and/or gearbox to be positioned within a pre-assembled cradle assembly and allows for easy access to the motor and gearbox for maintenance. The cradle assembly includes two end caps, a plurality of tube portions connecting the end caps and a plurality of motor mounts for supporting the motor and gearbox.

Hybrid and electric vehicles are enjoying increased popularity and acceptance due in large part to the cost of fuel for internal combustion engine vehicles. Electric axles for electric and hybrid vehicles can have an electric machine/motor with an output shaft that runs parallel to the axle of the motor vehicle. To support the electric machine/motor and a transmission, a cradle assembly is needed to provide sufficient strength and stiffness to support the weight of machine/motor and transmission and the load of the vehicle transferred to the cradle assembly.
The cradle assembly disclosed below provides the necessary strength and support for the electric axle assembly while allowing access to the electric machine/motor and/or transmission for maintenance and repair without requiring the disassembly of the cradle assembly itself.

Cradle assembly for an electric axle assembly


An exemplary multi-monitor method includes controlling a first position determining unit and a second position determining unit of a corresponding monitor to emit signals. Then, the method determines which of the first position determining unit and the second position determining unit of the corresponding monitor has received the signal. The method next determines the position of the corresponding monitor among all of the monitors, and determines a new coordinate range of the corresponding monitor in the Cartesian coordinate system. Then the method determines which point has been touched in response to a user touch on the touch unit, and determines the new set of coordinates in the Cartesian coordinate system corresponding to the touched point. Next, the method determines one or more touched monitors, and controls the determined one or more touched monitors to execute the corresponding operation.

1. Technical Field
The present disclosure relates to multi-monitor systems and, more particularly, to a multi-monitor system capable of controlling the operation of a number of monitors in response to user touch, and a multi-monitor controlling method.
2. Description of Related Art
In a conventional multi-monitor system, two or more monitors are connected to one host so that the monitors are used as one large monitor. However, in the conventional multi-monitor system, a mouse is connected to the host. The mouse must be moved frequently to operate a corresponding monitor. The need to frequently move the mouse is inconvenient to the user. Therefore, it is desired to be a multi-monitor method to resolve the above problem.

Multi-monitor system and multi-monitor controlling method

Hong Fu Jin Precision Industry (WuHan) Co., Ltd.