Column-stabilized floating structures having a deck and a plurality of vertical buoyant caissons bridged together in distantly spaced relation by a plurality of open frame horizontal truss pontoon members and vertical truss columns at a lower end. The buoyancy of the caissons is selectively adjusted by means of ballast control. Water is selectively pumped into or out of keel tanks at the bottom of the truss structure such that the water mass and weight is adjustably tuned to raise or lower the center of gravity of the entire mass of the floating structure relative to its center of buoyancy. By tuning and positioning the center of gravity relative to the center of buoyancy, the floating structures can be tuned according to ballast and other variable or fixed loads including the deck payloads and to compensate for different operational, environmental, and survival conditions and towing and installation stages.

1. Field of the Invention
This invention relates generally to apparatus and installation methods for semi-submersible vessels, tension leg platforms and floating structures, and more particularly to column-stabilized floating structures having a plurality of vertical spaced apart buoyant caissons connected together by an open framework of truss-pontoon members and the method of installation and operation.
2. Brief Description of the Prior Art
The word truss, as used herein, refers to a welded or bolted open frame structure formed of slender tubular members. Horizontally disposed truss structures are known as xe2x80x9ctruss pontoonsxe2x80x9d, and vertically disposed truss structures are known as xe2x80x9ctruss columnsxe2x80x9d. The truss bridges between three or more buoyant vertical column structures that stabilize a semi-submersible vessel at the water surface when floating with respect to wind, wave, current and other horizontal loads.
Semi-submersible type offshore floating vessels were first introduced in 1950. Ever since their introduction, the semi-submersible vessels have been found to be economical for offshore drilling, production, transport, pipe laying and other shallow and deep-water applications. SPAR, TLP, FPSO""s and barges are other types of floating vessels that can be used for similar applications as semi-submersible vessels. Decks, columns and pontoons are the typical components of a semi-submersible vessel. Semi-submersible vessels comprise three or more large size vertical columns that are spatially separated and connected by large size horizontal pontoons between the columns at their bottoms. The deck structure is above the water with a sufficient air gap between the still water level to the bottom of the deck to allow waves to pass across the columns without impacting on the bottom of the deck. The center of gravity of the entire semi-submersible vessel is generally high due to large deck loads above the water. The deck of a semi-submersible vessel is a very conventional structure made of boxes, trusses and or girders fixed on top of the columns. The modern semi-submersible vessel uses pleated box structures for the deck with bulkheads, frames and stiffeners, in which case, the deck is called the upper hull. The column is a shell type plated structure with stringers, bulkheads and stiffeners.
The pontoon is also a shell type plated structure with frames, bulkheads and stiffeners. The pontoon and columns have compartments for ballast, voids and storages. The column and the pontoon together is called lower hull. The lower hull provides the necessary buoyancy to the semi-submersible to take the functional structural, equipment and live loads. The vertical columns are sized and located apart to provide water plane area and to provide the area moment of inertia. Since the center of the gravity of the semi-submersible vessel is much higher above the center of buoyancy of the lower hull, the acting moments due the wind, wave and currents horizontal forces are restored by the semi column water plane areas and moment of inertia of the water plane areas. The size and the spacing of the columns provide these two stability parameters. Thus the semi-submersible design is called a xe2x80x9ccolumn-stabilized floating unitxe2x80x9d. The semi-submersible vessel is a free-floating vessel compliant to wind, waves and currents. It is moored and/or dynamically positioned by powered-thrusters.
Semi-submersible vessels have good sea-keeping behavior with respect to waves and have large deck areas and carrying capacity. However the current and wind resistance results in higher positioning. A semi-submersible vessel is sensitive to the variable deck loads and its stability is limited. Semi-submersibles have large heave motions (vertical oscillation of the floating vessel) when subjected to waves and consequently a dry-tree oil production system is not feasible. In the case of a drilling semi-submersible, downtime of drilling may occur when the heave motion is not feasible for drilling.
Thomas, U.S. Pat. No. 6,024,040 discloses a mobile jack-up platform converted to a semi-submersible offshore platform. The platform includes single submerged hollow lower base at the bottom end, partially submerged elongate vertical buoyant connecting legs extending upwardly therefrom and passing through an upper barge (jack-up platform) above the level of the sea. The hollow base has a square, rectangular or triangular configuration and is filled with seawater to form the ballast for the entire platform and may include interior reservoirs in which hydrocarbons are stored. A central opening or passage in the center of the base reduces the resistive surface size of the base in the water during vertical movements of the platform. The vertical connecting legs have a hollow cylindrical upper portion with a bottom wall forming a buoyancy tank, and a lower portion formed of open frame lattice construction. The respective lengths of the hollow cylindrical buoyant upper portion and the lattice-work lower portions are dimensioned relative to one another so that a pressure force exerted by the sea on the upper portions substantially compensates for an acceleration force exerted on the base by the action of the seawater surrounding the base over a usual swell period range of the sea. The platform in operation imitates a semi-submersible, which can retract the legs with respect to the upper barge (jack-up platform). With the legs fully retracted above the upper barge and its single lower base closely adjacent to the bottom of the upper barge, the platform can be floated and transported to another location.
The tension leg platform (TLP) was introduced as a new concept based on a semi-submersible design with a deck, columns or caissons and pontoons. The vertical heave of the tension leg platform (TLP) is reduced by tendons attached between the lower hull and the seabed. The tendons are always maintained in tension with the excess buoyancy designed into the lower hull. The application of a TLP is limited to water depths of say 5000 ft., over which, the tendons"" vertical neutral period enters the energetic portion of the wave heave forces of the TLP. Secondly, the size of the lower hull increases the economics to always maintain tension in the tendons with respect to the increased water depth.
The SPAR is currently used in relatively large offshore oil production applications. The SPAR is a single large vertical column or caisson buoyant structure. Unlike the semi-submersibles, the stability of the vessel is not dependent on the water plane area and the moment of inertia of the water plane area. The stability of the SPAR is provided by lowering the center of gravity of the vessel below its center of buoyancy. Thus, the vertical buoyancy acts upwards above the center of gravity, and total weight acts at the center of gravity below the center of buoyancy. The SPAR has excellent heave performance like a TLP and dry-tree production systems like a fixed offshore platform are used. However, the SPAR has limitation in deck area and payload due to the size of the SPAR. The SPAR is fabricated and transported horizontally and up-ended vertically at the operating location. Then the fully equipped deck is installed on top of the SPAR vertical column. Thus the transportation and installation cost and risk are significantly increases with water depth and deck payload applications. The SPAR becomes uneconomical for larger production as a self-contained drilling unit in deepwater. Ultra deepwater posses problems to the riser tensioning system and limits the applicability of the SPAR. Mooring of the SPAR in deepwater over 5000 ft also posses a serious problem. The effectiveness of the mooring and handling are reduced in such water depth. The SPAR also poses vortex shedding vibration problems when the vertical single column hull becomes extremely slender.
Horton, U.S. Pat. No. 5,558,467 discloses a spar-type deep water offshore floating apparatus, for use in oil drilling and production in which an upper buoyant hull of prismatic shape is provided with a passage longitudinally extending through the hull in which risers run down to the sea floor, the bottom of the hull being located at a selected depth dependent upon the wind, wave, and current environment at the well site, which significantly reduces the wave forces acting on the bottom of the hull, a frame structure connected to the hull bottom and extending downwardly and comprising a plurality of vertically arranged bays defined by vertically spaced horizontal water entrapment plates and providing open windows around the periphery of the frame structure, the windows providing transparency to ocean currents and to wave motion in a horizontal direction to reduce drag, the vertical space between the plates corresponding to the width of the bay window, the frame structure being below significant wave action whereby wave action thereat does not contribute to heave motion of the apparatus but inhibits heave motion, the frame structure serving to modify the natural period and stability of the apparatus to minimize heave, pitch, and roll motions of the apparatus. A keel assembly at the bottom of the frame structure has ballast chambers for enabling the apparatus to float horizontally and for stabilization of the apparatus against tilting in vertical position, and taut anchor lines connected to the apparatus at allocation of relatively little cyclic movement of the apparatus, the said lines being connected to suitable anchors. Horton, III, U.S. Pat. No. 5,722,797 discloses a spar-type buoyant floating caisson for offshore drilling and production that includes means for increasing the natural period of the caisson and reducing heave, pitch, and roll without increasing the overall length of the caisson. The floating caisson has a center well through which drilling and/or production risers pass and one or, more circular plates extend radially from the caisson below the water surface. The circular plates provide additional mass and resistance to environmentally induced motions and thus increases the natural period of the caisson beyond the periods of maximum wave energy, which allows the caisson to be designed with a shallower draft than a caisson without the plates that would normally be used in deep water.
Blevins et al, U.S. Pat. No. 6,206,614 discloses a spar-type floating offshore drilling/producing structure that is formed from a plurality of closely spaced vertically oriented buoyant columns on which one or more modules or decks may be placed to support process equipment, a drilling rig, utilities, and accommodations for personnel. The columns are held in the closely spaced relationship by a plurality of horizontal plates spaced along the length of the columns and vertical plates located near the bottom of the columns and near the top of the columns. The columns have a smaller water plane area than the horizontal plates. The structure includes fixed ballast, an oil storage area, and voids and variable ballast for offsetting the lighter weight of the stored oil.
A long-felt need exists for an improved offshore platform that can economically support a deck in deepwater and respond efficiently to the environmental forces in ultra deepwater. An ideal vessel for a deepwater application requires: a large deck area, large carrying capacity, less heave response to waves, excellent wind and wave current motion behavior so that a dry-tree production system is feasible, less station keeping forces, and finally, less construction costs.
The present invention is distinguished over the prior art in general, and these patents in particular by a new concept for offshore drilling, production, transportation and other applications with the above noted key principle requirements for a feasible economical system. The present structure herein referred to as xe2x80x9ccolumn-stabilized floating structures with truss pontoonsxe2x80x9d has a deck and a plurality of vertically oriented buoyant columns or caissons bridged together in distantly spaced relation a plurality of open frame horizontal truss pontoon members and vertical truss columns at a lower end. Each of the buoyant caissons has a submerged portion and a non-submerged portion with one or more buoyant tanks or chambers enclosed by bulkheads or entrapped inside the columns or caissons, the buoyancy of which can be selectively adjusted by means of ballast control. Water is selectively pumped into or out of keel tanks at the bottom of the truss structure such that the water mass and weight is adjustably tuned to raise or lower the center of gravity of the entire mass of the floating structure. By tuning and positioning the center of gravity relative to the center of buoyancy, the present floating structures can be tuned according to ballast and other variable or fixed loads including the deck payloads and compensate for different operational, environmental, and survival conditions and towing and installation stages of the floating structure.
The present column-stabilized floating structures having buoyant caissons with horizontal truss pontoons and vertical truss columns provides a large deck area and deck payload like a conventional semi-submersible, and low heave response to waves and environmental forces like a SPAR vessel. Thus, the present invention has all of the advantages of a semi-submersible and a SPAR vessel. However, unlike the SPAR vessel concept, which requires the center of gravity to be below the center of buoyancy in order to achieve the required stability when floating, the center of gravity of the present structure can be selectively positioned above the center of buoyancy and is stabilized by distantly positioned buoyant caissons, similar to a semi-submersible vessel. By the fundamentals of the principles of the stability criteria, the present invention does not belong to the family of SPAR structures but stabilized by its multiple columns water plane area and the moment of inertia of the water plane area. Spacing the columns well apart increases the moment-of-inertia of the water-plane-area of the water piercing columns. Thus the present invention is called a column stabilized floating structure with truss pontoons.
It is therefore an object of the present invention to provide an offshore column-stabilized floating structure that can economically support a large deck area in deepwater and responds efficiently to environmental forces in ultra-deepwater.
It is another object of this invention to provide a column-stabilized floating structure that has a large carrying capacity and excellent heave response to wind and waves thereby making a dry-tree oil production system economically feasible.
Another object of this invention is to provide a column-stabilized floating structure that is stabilized by buoyant columns or caissons and allows the center of gravity to be above the center of buoyancy in order to achieve stability when floating.
Another object of this invention is to provide a column-stabilized floating structure wherein the center of gravity and weight of the structure can be adjustably tuned as desired to achieve a proper height by pumping water in and out of the keel tanks of the structure.
Another object of this invention is to provide a column-stabilized floating structure that has a low lower hull weight, reduced fabrication cost, and large payload to hull displacement ratio.
Another object of this invention is to provide a column-stabilized floating structure that has low mooring or DP positioning forces, less plated shell type structure and fatigue sensitive spots, and greater residual structural strength than conventional semi-submersible vessels.
A further object of this invention is to provide a column-stabilized floating structure that is insensitive to variable deck loads, allows all equipment and consumable fluids to be located on the deck, and has no risk of damage by dropped objects.
A still further object of this invention is to provide a column-stabilized floating structure that reduces engineering cost, reduces fatigue damages, is economical to manufacture, operate, and maintain.
Other objects of the invention will become apparent from time to time throughout the specification and claims as hereinafter related.
The above noted objects and other objects of the invention are accomplished by a column-stabilized floating structure having a deck and a plurality of vertically oriented buoyant caissons bridged together in distantly spaced relation by a plurality of open frame horizontal truss pontoon members and vertical truss columns at a lower end. Each of the buoyant caissons has a submerged portion and a non-submerged portion with one or more buoyant tanks or chambers enclosed by bulkheads or entrapped inside the columns or caissons, the buoyancy of which can be selectively adjusted by means of ballast control. Water is selectively pumped into or out of the keel tanks at the bottom of the truss structure such that the water mass and weight is adjustably tuned to raise or lower the center of gravity of the entire mass of the floating structure. By tuning and positioning the center of gravity relative to the center of buoyancy, the present floating structures can be tuned according to ballast and other variable or fixed loads including the deck payloads and compensate for different operational, environmental, and survival conditions and towing and installation stages of the vessel.

Column-stabilized floating structures with truss pontoons


A novel power conservation scheme is provided for conserving power in client terminals by using a proxy server. The client terminal may have a secondary communication interface for short range communications and a primary communication interface for long range communications with an access point for a wireless network. To conserve power, the client terminal may power down its primary communication interface without informing the access point. The access point assumes the primary communication interface is still active. Prior to shutting off its primary communication interface, the client terminal may assign an external proxy device to act as a proxy and monitor its data channel with the access point. The proxy device monitors the data channel(s) for the client terminal via a primary communication interface. If the proxy device detects a data message for the client terminal, it forwards the data message to the client terminal via a secondary communication interface.

1. Field
Various features pertain to wireless communication systems. At least one feature pertains to devices and methods for facilitating power conservation in a wireless client terminal by using a proxy device to receive and forward data to the client terminal via a secondary wireless interface.
2. Background
Client terminals, such as laptop computers, personal digital assistant devices, mobile or cellular phones, or any other device with a processor, that communicate with other devices through wireless signals are becoming increasingly popular. Client terminals are typically battery-powered and the amount of power a battery can provide is limited. With consumers using power-intensive applications that run on the client terminals, particularly entertainment media and imaging applications, conserving battery power is important.
Client terminals may operate in one of several modes depending on the communication standard being implemented in a wireless network. For example, Wideband Code Division Multiple Access (WCDMA) includes various types of "connected" modes (e.g., Radio Resource Control states—Cell Dedicated Channel DCH, Cell Forward Access Channel FACH, Cell Paging Channel PCH, and UMTS Terrestrial Radio Access Network (UTRAN) Registration Area (URA) Paging Channel PCH) and an idle mode. Typically, in the connected mode a client terminal may be fully or partially connected to a wireless network to receive and/or transmit data (e.g., voice or data calls or sessions) with one or more access points (e.g., base stations, Node B, femto cell, etc.). In connected mode, conventional client terminals may consume a significant amount of power to sustain circuitry needed to monitor a data channel and/or maintain a connection with the wireless network. In idle mode, the client terminal may monitor control channels, such as a paging channel (PCH) for paging messages. Such paging messages may include messages that alert the client terminal to the occurrence of an incoming data message (e.g., a data packet or segment) and control/overhead messages that carry system information and other information for the client terminal. While power consumption in idle mode is substantially less than in connected mode, there is a longer delay in powering the circuitry and re-establishing a wireless network connection from idle mode if a message is received. This delay is relatively long (e.g., in comparison to delays in connected mode) and would impact the user experience noticeably (e.g., noticeable delay in receiving a message, a dropped message, increased retries from sender, etc.). Thus, maintaining the wireless terminal in a connected mode results in faster power depletion or consumption but switching to idle mode results in a noticeable lag or delay.
Therefore, there is a need for a solution that reduces power consumption of a client terminal while avoiding noticeable delays typically associated with power-conserving modes operation.

Power conservation in wireless client terminals using proxy device


A method of constructing a globally unique locator for a network entity attached to a network. A path between a core network and the network entity via intermediate network entities is determined by applying policy routing. The globally unique locator for the network entity is constructed by concatenating elements identifying the locations of the network entity, the intermediate network entities and a node in the core network. The identifying elements are obtained from Attachment Registers associated with each of the network entity, the intermediate network entities and the node in the core network.

Name-address management generally includes issues such as name-to-address resolution and name-address registration. Name-to-address resolution is a procedure by which a "name" of a network resource, e.g., a network node, is resolved or translated into a routable network address, i.e., a location in the network topology. Name-address registration is the corresponding registration procedure by which the name and the assigned network address of the resource are registered in the network. The name of the resource is normally known to users and typically also stays the same over relatively long periods of time.
Traditional networking architectures for the Public Switched Telephone Network (PSTN) or Internet solve the problem of connecting terminals or hosts (which can be considered as "boxes") in order to support a specific application such as telephony or WWW. To this end, traditional naming and addressing schemes employ host or network centric identities such as E.164 numbers for telephony, or Internet Protocol (IP) addresses and Uniform Resource Locators (URLs) for the Internet. However, the end user is typically interested in reaching a destination object that sits behind or within the box, such as a human being or a file, rather than communicating with the box itself. As the destination objects move to new boxes, box or network dependent identities of these objects must be updated. For example, it is an everyday experience that a person cannot be reached at the phone number registered in some semi-static directory because that directory does not track the phone that she is momentarily close to. Or a web link is broken because the target data object has been moved to another location. To address this class of problems, box-independent addressing schemes such as Telephone Number Mapping (ENUM), Session Initiation Protocol (SIP) names, and Uniform Resource Identifiers (URIs) have been developed. Using a presence or mobility mechanism, box-independent names of a destination object can then be mapped to the address of the box where the destination object is currently located. Likewise, inventory systems are used to track specific objects to a specific location.
In addition, it is quite common that mobile objects travel within or in association with other mobile objects. For example, a person may travel in association with a mobile phone that may be associated with a personal area network that may be associated with a vehicular network. Consider, for example, a person using a laptop with a Wireless Local Area Network (WLAN) who is working onboard a train that has Internet connectivity. To reach this person, three separate mobility-related functions must be used: 1) a presence function that binds the person to the laptop, 2) a mobility function that binds the laptop to a specific IP address on the train, and 3) a mobility function for the train.
Another example is a piece of merchandise that travels in a container that travels on a ship. For this case a separate inventory system would be used with little in common with the functions for personal and host mobility described in the previous example.
The Domain Name System (DNS) stores and provides information associated with domain names in a distributed database in networks such as the Internet. The DNS associates domain names with many types of information, but most importantly it provides the IP address for a given domain name. DNS makes it possible to associate easy-to-remember domain names (such as ericsson.com) with hard-to-remember IP addresses. DNS is suitable for resources that rarely change their location, but is not adapted for mobility. RFC 2136 describes "Dynamic Updates in the Domain Name System" in the hope of providing better support for rapid updates of DNS, but it is still far from being suitable for keeping track of roaming resources such as mobile phones and their users.
When routing protocols for the Internet and other fixed networks were initially created, hosts were not expected to move around. Therefore, hosts are usually named by their point of attachment to the network, e.g., IP addresses. Examples of such routing protocols include RIP, IS-IS, OSPF, BGP and PNNI. They are all well established technologies but have limited support for mobility and have convergence problems when network topologies change rapidly.
Traditionally, applications use IP-addresses in a way that does not allow them to change during an on-going session. To allow hosts to move without changing their IP-addresses (at least from an application perspective) mobility solutions in IP networks, such as Manet (ad hoc networks), Network Mobility (NEMO), and Mobile IP have been developed. But these are fairly complex solutions since they adapt a technology intended for fixed networks to new mobility requirements. In addition, there is a multitude of inventory systems for various classes of objects, such as library books or pieces of merchandise. Each system and mechanism is optimized for its specific class of objects. For example, mobile IP is optimized for mobility of IP boxes, SIP supports personal mobility, inventory systems handle mobility of goods, etc. Due to the diverse technologies employed in this field, it is hard to achieve synergies between the systems and mechanisms used for the different classes of mobile objects.
The Host Identity Protocol (HIP) provides a method of separating the end-point identifier and locator roles of IP addresses. It introduces a new Host Identity (HI) name space based on public keys. The public keys are typically self-generated. The HIP separation can be used to provide end-to-end connectivity over different locator domains. Even still, routing protocols under development cater to mobility of individual hosts (nodes) but do not adequately address the problems relating to mobile networks (MNs). A mobile network includes a group of many mobile hosts or other objects that move together as a group. Mobile network examples include networks located in any type of mobile vehicle, e.g., in a train, airplane, bus, ship, subway, etc., but are not limited to vehicles. All that is required is that the group of mobile objects, hosts and routers move substantially together at substantially the same time. Also, a communication satellite carrying a router is another example of a mobile network that dynamically attaches to ground stations, other communication satellites, and hosts or mobile phones. A particular mobility problem associated with mobile networks is a potentially huge number of registration or other location updates that need to be signalled and processed whenever the mobile network changes location. Such a move may cause an "update storm."
Consider for example a Public Land Mobile Network (PLMN) type of system like GSM and 3G cellular networks. Mobile host name resolution is handled via a Home Location Register (HLR) and the Visited Location Register (VLR). When a mobile host is called, a phone number (MS-ISDN) is resolved via the VLR and HLR into a corresponding E.164 address that allows the call to be routed to the mobile host, if the mobile host with the MS-ISDN has registered its current location area with the VLR. Local mechanisms are used to route the call to the specific cell in the location area in which the mobile host is currently located.
The HLR and VLR have overall good performance and security support regarding name resolution in cellular systems. But they are closely linked to the E.164 address structure and as such do not provide an open architecture for other and/or arbitrary name and address spaces. Moreover, this approach to registering a host with a centralized location register like the HLR/VLR does not function well with mobile networks. The problem is particularly acute when a large mobile network with many subnetworks or hosts roams and requires registration update signalling for every one of its subnetworks and/or hosts—a good example of an "update storm" mentioned above.
In dynamic DNS, when such a mobile network roams, each host in the mobile network must have its DNS record updated. For that situation, mobile IP requires that all home agents having hosts in the mobile network be updated. RFC 3963 describes the IETF Network Mobility (NEMO) basic support protocol which enables mobile networks to attach to different points in the Internet. The protocol is an extension of mobile IPv6 and allows session continuity for every node in the Mobile Network as the network moves. It also allows every node in the Mobile Network to be reachable while mobile around. But NEMO's distributed solution suffers from what is called "pinball routing," where all internetwork traffic must be routed between every mobility agent that has an associated mobile node or network in the path towards the destination host. As a result, tunnelling overhead accumulates per radio hop, and there are potential latency problems when several mobility agents are located at different continents, i.e., x-ogonal routing instead of triangular routing.
As illustrated in FIG. 1, it is known to provide a set of Attachment Registers 14 in a network 12, where the attachment registers perform one or more network functions. A primary function of the attachment registers is to store a locator point of an object reachable from a network. For example, Attachment Register A (ARA) stores a locator point of network A 16. Information is stored in the attachment registers that establishes one or more logical links between attachment registers. In the example shown in FIG. 1, the arrows between the attachment registers 14 indicate that information stored in attachment register ARD links to attachment register ARC. Similarly, information in ARC links to ARB, and information in ARB links to ARA. The information in these attachment registers creates a path to host D, and thus allows easy construction of a global locator to host D. By using attachment registers, a global locator can be constructed to a host attached to a network even where the network or an intermediate network changes location, without all hosts attached to the network having to change their address. So if, for example, Network C (Nc) changes it's location, ARc is updated with the new location of Nc. The AR for the other networks and network entities need not change.
The constructed locator consists of a list of nodes or links that describes a specific path between the Highest Level Network (or core network) and host D. In the general case when networks below the Highest Level Network are interconnected in an arbitrary topology, there are many possible paths between the Highest Level Network and host D. The locator construction procedure must therefore be guided by a routing protocol operating between the Attachment Registers that selects the most optimal path. Examples of policy routing include finding a path with a specific bandwidth, packet loss or delay capability, or finding a path that traverses at least one network with a specific functionality, such as transcoding. An additional example of policy routing is to find a path that traverses networks that belong to a specific operator.

Policy routing in a communications network


A thermotropic chiral nematic liquid crystalline copolymer composition comprises ##STR1## wherein --NEM-- and --NEM'-- are each independently nematogenic units of formula ##STR2## wherein R is H or CH.sub.3, --Q-- is an alkylene radical having 1 to about 8 carbon atoms,. --X-- is --O--, --S--, or --CH.sub.2 --, --Y-- is ##STR3## --Z is --CN, --NO.sub.2 or --N.dbd.C.dbd.S, q and r are each independently 0 or 1; PA1 --X'-- is --O--, --S--, or --CH.sub.2 --, PA1 --Z' is an alkoxy, aralkoxy, alkylamino, or aralkylamino radical having 4 to about 12 carbon atoms and containing at least one asymmetric carbon atom, PA1 q' and r' are each independently 0 or 1; PA1 and wherein x is the mole fraction of chiral units and (y+y') is the total mole fraction of nematogenic units in said copolymer composition, and the ratio of x to (y+y') is from about 1:50 to 1:1.. wherein --CHI-- is a chiral unit of formula ##STR4## wherein R is H or CH.sub.3, --Q'-- is an alkylene radical having 1 to about 8 carbon atoms,. This compolymer composition is employed to form an optical device.

Liquid crystalline materials exhibiting the cholesteric mesophase have been proposed for use in a variety of optical applications, for example, notch filters, circular polarizing filters, selective reflectors, beam splitters, and beam apodizers. U.S. Pat. No. 3,711,181, for example, discloses an optical apparatus for modulating circular-polarized light that contains optically negative liquid crystal films.
It has been recognized that low molecular weight liquid crystal materials suffer limitations as to durability, effective temperature, mesophase stability, and amenability to processes for device fabrication. Polymeric liquid crystalline compositions have been proposed as potentially useful for overcoming such limitations. U.S. Pat. No. 4,293,435 discloses a liquid-crystalline cholesteric polymer phase that consists essentially of a copolymer of particular nematogenic and chiral acrylic ester monomers. In U.S. Pat. No. 4,410,570 is disclosed a liquid crystalline phase that contains a cyclic organopolysiloxane to which is chemically bonded at least one mesogenic group. Thermotropic cholesteric liquid crystalline glutamate copolymers consisting of chiral glutamate ester repeating units are disclosed in U.S. Pat. No. 4,743,675.
Many applications of chiral liquid crystalline polymeric materials in optical devices require polymers capable of forming both right- and left-handed helical structures. When a film of such a polymer is applied to a substrate, the helical structures must be capable of forming and maintaining the Grandjean texture, in which the helical axis is perpendicular to the substrate surface, to enable the selective reflection of circular-polarized light. An enantiomeric chiral pair of liquid crystalline polymers, whose individual structures are characterized as a right-handed and a left-handed helix, are thus capable of selectively reflecting fight-handed and left-handed circular-polarized light, respectively.
Especially useful for optical information storage applications are chiral nematic liquid crystalline polymers which form clear, transparent films that absorb no light in the visible region but do selectively reflect visible circular-polarized light. It is especially desirable for the application of these films as high efficiency polarizers that the reflected light be characterized by a broad half band width (HBW), as defined by the width of the spectral band measured at one-half of its maximum height. In devising flat panel displays, for example, where electrical power requirements should be kept as low as possible, chiral nematic liquid crystalline copolymers that form films whose half band widths in the visible region are substantially broader than those known in the art would be extremely useful.
In addition to the just described optical characteristics, it is necessary that the chiral nematic polymers be readily synthesized and that they have solubility properties which enable their processing into clear glassy thin films for use as optical devices. All of these requirements are met by the chiral nematic liquid crystalline copolymer compositions of the present invention.

Thermotropic chiral nematic liquid crystalline copolymers


An acoustic decoy round ejected by a launcher for flight above water from a ea-going vessel, impacts at the water surface to cause separation of a payload from a forward section of the round that is also separated from a flotation anchor tethered to the payload and fins which stabilized launched flight of the round prior to impact. The separated payload submerges from the flotation anchor at the water surface location to a tethered depth within the water from which a decoy signal is emitted.

This invention relates generally to acoustic devices launched from marine seagoing vessels or the like.
Expendable acoustic devices are presently available for deployment into seawater to counter torpedoes targeted at naval submarines. Currently, naval surface ships deploy torpedo decoys for submarines by manually activating and throwing the acoustic devices overboard. Such manual deployment of acoustic devices from surface vessels is often excessively time consuming, difficult to perform under high sea conditions, and often inadequate because of separation between the ship and the final deployment location of the acoustic device in the water.
It is therefore an important object of the present invention to provide for deployment of acoustic devices from sea-going surface vessels by use of missile launchers.
In accordance with the foregoing objective, it is a more particular object of the present invention to provide a transport round for launch of an acoustic type payload from existing mortar or rocket launchers on naval surface vessels, whereby safe, rapid and effective deployment of a torpedo decoy is achieved.

The United States of America as represented by the Secretary of the Navy


The present disclosure describes methods and systems for determining a flow of a combustible portion of vent gas delivered to an engine. The flow rate measurement may be performed by using the engine response to a relatively short (e.g. 1 to 5 s) interruption of the vent gas flow. A cross-correlation between RPM data of the engine and a reference signal corresponding to a state of a valve configured to interrupt the vent gas flow is determined, and a flow rate of the combustible portion of the vent gas delivered to the engine is determined from the maximum value of the cross-correlation.

In many industrial processes vent gases with combustible components are released to the atmosphere. These gases may contain undesired components such as benzene or methane. In the technology described by Malm in U.S. Pat. No. 8,113,181 entitled "Method and Apparatus for Capturing and Controlling Fugitive Gases", U.S. Pat. No. 8,235,029 entitled "Method and Apparatus for Processing Diluted Fugitive Gases", U.S. Pat. No. 8,382,469 entitled "Method and Apparatus for Utilising Fugitive Gases as a Supplementary Fuel Source", U.S. Pat. No. 8,851,054 entitled "Apparatus for Providing Diluted Fugitive Gases as a Fuel to an Engine", U.S. Pat. No. 8,978,627 entitled "Method and Apparatus for Capturing and Controlling Fugitive Gases", and U.S. Pat. No. 9,534,564 entitled "System and Method for Controlling a Flow of Vent Gases to a Natural Gas Engine", the entire contents of which are incorporated herein by reference, vent gases with gaseous combustible components are added to the intake air of an internal combustion engine. This reduces the need for engine fuel and efficiently eliminates these components from the vent gases. In the case of methane, which is a potent greenhouse gas, there is a need to measure the flow rate to determine the benefit of the technology. Flow measurement is done either using a calibrated flow meter or calculation using the pressure drop due to a restriction such as an orifice or valve.
Additional, alternative, and/or improved systems and methods that enable for determining the flow of the combustible portion of vent gas remain highly desirable.
It will be noted that throughout the appended drawings, like features are identified by like reference numerals.

Measurement of flow of vent gas with combustibles


A radio-frequency identification (RFID) system including an RFID tag and an RFID-enabled object. The RFID tag may include a pair of antennas and an RFID circuit. The antennas receive activation energy from a reader, and the RFID circuit modulates tag energy when activation energy is received by one of the antennas. The tag may also include a transmission line for operatively coupling the RFID circuit to the antennas. A first one of the antennas may receive activation energy which, in turn, may be radiated by a second one of the antennas. The second antenna may also receive tag energy radiated by an antenna of another RFID tag. The received tag energy may then be radiated by the first antenna. Accordingly, when a plurality of the RFID tags are positioned in sequence in operative proximity with each other, activation energy may be propagated through the sequence from one RFID tag to another in one direction, and tag energy may be propagated through the sequence from one RFID tag to another in another direction.

The present invention relates to radio-frequency identification (RFID) systems, including RFID tags and readers. The invention also relates to RFID apparatus and methodology that enables a plurality of cartons, such as stacked on a pallet, to be read-even the innermost cartons in the stack and even in the presence of RF intolerant material such as liquids and metals.
Automatic identification is the broad term applying to a host of technologies that are used to help machines identify objects. Automatic identification is often coupled with automatic data capture. More specifically, companies want to identify items, to capture information about the items, and to load the information into a computer with minimal human labor.
One type of automatic identification technology is radio-frequency identification (RFID). RFID is a generic term for technologies that use radio waves to automatically identify individual items. There are several conventional methods of identifying objects using RFID, the most common of which is to store a serial number (and other information if desired) that identifies a product on a microchip that is attached to an antenna. The chip and the antenna together are called an RFID transponder or an RFID tag. The antenna enables the chip to receive commands from and to transmit identification information to a reader. The reader converts the radio waves returned from the RFID tag into a form that can then be utilized by a computer.
Accordingly, a conventional RFID system consists of a tag (including a chip with an antenna) and a reader (sometimes call an interrogator) with an antenna. The reader sends out electromagnetic waves that form a magnetic field when coupled with the antenna on the RFID tag. A passive RFID tag draws power from this magnetic field and uses the power to drive or activate the chip. The chip then modulates the waves that are sent back to the reader, and the reader converts the returned waves into digital data.
There are generally two types of RFID tags: active and passive. An active RFID tag utilizes a battery to power the chip and to transmit a signal to a reader (similar to a cell phone transmitting signals). A passive tag does not have a battery but rather is powered by the electromagnetic waves that induce a current in the antenna of the tag. A semi-passive tag uses a battery to power the chip but communicates using electromagnetic waves from the reader.
Similar to an audio radio tuning in to different frequencies, RFID tags and readers are tuned to the same frequency to communicate. RFID systems use many different frequencies, but the most common frequency ranges utilized in RFID systems are low-frequency (about 125 KHz), high-frequency (13.56 MHz), and ultra-high frequency or UHF (about 900 MHz). Microwaves, which have a frequency of about 2.45 GHz, are also used in some applications.
The distance at which an RFID tag can be read is known as the read range. The read range of a passive tag depends on a number of factors: the frequency of operation, the power of the reader, and interference from metal objects or other RF devices. In general, low-frequency tags have a read range of about one foot; high-frequency tags have a read range of about three feet; and UHF tags have a read range of about 20 feet. Where longer read ranges are needed, an active tag with a read range of 300 feet or more can be used.
One of the desired applications of RFID tags is to track goods in a supply chain, particularly at high volumes such as a plurality of cartons stacked on a pallet. One of the inherent difficulties in this application is ensuring that all of the RFID tags associated with all of the cartons are read. For example, if the stack of cartons is a five-by-five layer stacked five layers high (i.e., 125 cartons total), then the user would want to ensure that all 125 RFID tags are read, even those tags mounted to cartons located in the center of the stack. This effectiveness of this operation may be aggravated by the presence of material that is not conducive to RF reading.
More specifically, radio waves are able to travel through most non-metallic materials, so that RFID tags can be embedded in packaging or encased in protective plastic for weather-proofing and durability while still being readable. However, radio waves reflect off metal and are absorbed by water and complex organic molecules such as fat at higher frequencies. These characteristics make tracking metal products or those with high water content problematic. In addition, reading a stack of cartons with RFID tags, particularly cartons located in the center of the stack or cartons that may contain metals or liquids, is also problematic.
In view of the foregoing, there is a need in the art for RFID technology that enables all of the RFID tags in a stack of cartons to be read. The present invention satisfies this need.

This patent is subject to a terminal disclaimer.

RFID tags for enabling batch reading of stacks of cartons


A technique for finding an optimal threshold for binarization of a gray scale image employs fuzzy reasoning. A triangular membership function is employed which is dependent on the degree to which the pixels in the image belong to either the foreground class or the background class. Use of a simplified linear fuzzy entropy factor function facilitates short execution times and use of membership values between 0.0 and 1.0 for improved accuracy. To improve accuracy further, the membership function employs lower and upper bound gray level limits that can vary from image to image and are selected to be equal to the minimum and the maximum gray levels, respectively, that are present in the image to be converted. To identify the optimal binarization threshold, an iterative process is employed in which different possible thresholds are tested and the one providing the minimum fuzzy entropy measure is selected.

1. Field of the Invention
The present invention relates in general to a method and system for converting gray scale images to binary images which employs fuzzy reasoning to calculate an optimal binarization threshold value.
2. Description of the Background Art
Conversion of gray-scale digital images to binary images is of special interest because an image in binary format can be processed with very fast logical (Boolean) operators by assigning a binary value to each of the image's pixels. A binary one value indicates that the pixel belongs to the image foreground, which may represent an object in the image, while a binary zero value indicates that the pixel is darker and belongs to the image's background. Since most image display systems and software employ gray-scale images of 8 or more bits per pixel, the binarization of these images usually takes 2 extreme gray tones, black and white, which are ordinarily represented by 0 or 255, respectively, in an 8-bit gray-scale display environment.
Image thresholding is the simplest image segmentation approach for converting a gray-scale image to a binary image. It is actually a pattern classification procedure in which only one input feature is involved, this being the pixel intensity value. Usually a binary image is obtained from an 8-bit gray-scale image by thresholding the image and assigning either the low binary value (0) or the high (255) value to all gray levels based on the chosen threshold. Obviously, the threshold that is chosen has a critical importance since it controls the binary-based pattern classification that is obtained from the gray-scale image. The key issue is to choose an optimal threshold so that the number of misclassified image pixels is kept as low as possible. Since images can differ substantially from one another depending on the objects contained therein, the optimal threshold value can vary considerably from one image to the next. Thus, merely selecting a threshold value that is, for example, set at the average pixel intensity value for the gray-scale image will probably not provide the optimal threshold. If the threshold is selected incorrectly, substantial image information will likely be lost in the conversion to binary.
Numerous techniques have been employed to address the foregoing issue. The most accurate of these are non-interactive techniques that do not require selection of any process parameters to identify the optimal threshold. Such techniques automatically select the appropriate threshold based on an analysis of each image to be converted. An example of such a technique is disclosed by N. Otsu in A Threshold Selection Method From Gray-Level Histograms. IEEE Transaction on Systems, Man, and Cybernetics, 9(1):62-66, (1979) (hereinafter referred to as the Otsu method). In the Otsu method, the optimal threshold is determined by minimizing the two variance classes; total variance and in-class variance. In other words, the means/averages of the two classes (background and foreground) should be as well separated as possible and the variances (standard deviation) in both classes should be as small as possible. The Otsu method is basically based on selecting the lowest point between the two classes.
One particularly promising non-interactive approach is to employ fuzzy reasoning to determine the optimal threshold for binarization. Fuzzy reasoning is a logical reasoning technique that attempts to mimic more accurately how the human brain reasons. Under the fuzzy reasoning approach, a logic problem becomes more than deciding whether to assign a binary one or zero to a particular bit, pixel or parameter. Fuzzy reasoning goes one step further and recognizes that there is information contained in the degree to which a given value possesses a particular characteristic. For example, there is much less certainty that a particular pixel is in the background or foreground of the image if the pixel is very near a selected intensity threshold than if the value were far below or above the threshold. In a fuzzy reasoning approach, a multiple pixel digital image is defined as an array of fuzzy singletons, each having a membership value somewhere between 0.0 and 1.0 that denotes its degree of possessing some property (e.g., brightness, darkness, edginess, blurredness, texture etc.). For image binarization, the membership function is defined in terms of the degree a pixel having a particular gray level value in the image belongs to one of the two binary classes, background and foreground.
Once the membership function is formed, the function can be employed to determine the optimum threshold value that defines the boundary between background and foreground gray levels. This is accomplished by identifying the threshold value which results in the membership function providing the minimum fuzzy entropy for the image. The concept of fuzzy entropy is generally defined in information theory as a measure of information. In the context of fuzzy reasoning, the entropy is a measure of the degree of fuzziness. Thus, in the image binarization application, the goal is to select a threshold value that results in the minimum fuzziness or uncertainty.
An example of the use of fuzzy reasoning in image binarization is the method disclosed by Huang and Wang in Image Thresholding by Minimizing the Measures of Fuzziness, Pattern Recognition, Vol. 28, No. 1, pp 41-51 (1995) (hereinafter referred to as the Huang-Wang method). In the Huang-Wang method, a triangular membership function for the foreground and background classes is employed in which the graph of the function appears as two adjacent triangles that join at a selected threshold value. The peak values of the triangles occur at the average pixel intensity level for each class, where the membership value is 1.0. To identify the optimal threshold, an iterative trial and error technique is employed to identify the threshold that results in the minimum fuzzy entropy for the membership function. Shannon's entropy function, which is a logarithmic function in the shape of a parabola, is used as an entropy factor or cost function to calculate the entropy measure for a selected threshold. The threshold value that results in the minimum fuzzy entropy is then selected as the optimal threshold for binarization of the image.
Although the Huang-Wang method is fairly accurate and selects image thresholds that in general result in preservation of more image information than more conventional techniques, this increased conversion accuracy comes at the expense of substantially more computational power and execution time. For example, In tests comparing the Huang-Wang method to the Otsu method, the Huang-Wang typically required approximately 3 times the execution time than that of the Otsu method. The extended execution time is primarily due to the logarithmic nature of Shannon's entropy function which complicates the necessary calculations. In addition, use of Shannon's function restricts the values of the membership function to a range of 0.5 to 1.0, which limits accuracy. The limited range is necessary because the parabolic shape of Shannon's function has increasing values between membership values of 0.0 and 0.5, and decreasing values between membership values of 0.5 and 1.0. However, because the cost or entropy should decrease as the membership function value increases (as the fuzziness becomes smaller), the membership values below 0.5 cannot be employed when Shannon's function is selected as the entropy measure function. As a result of the foregoing, there is a need for a fuzzy reasoning based binarization technique that can operate effectively with reduced execution times.

Optimal binarization of gray-scaled digital images via fuzzy reasoning

United States of America as represented by the Administrator of the National Aeronautics and Space Administration


The invention relates to a method for packing and filtering geophysical events read on multi-domain records, with distribution of said events into families each having a particular geophysical significance: iso-offset or iso-incidence angle data cubes, elastic parameter cubes derived from a joint stratigraphic inversion and the like, in order to extract therefrom data on the type of subsoil, using a multivariate statistical technique. The method essentially comprises forming, by combination of geophysical variables, synthetic variables much fewer in number, which are obtained by constructing an orthogonal vector base in each of the sets of analysis constituted by the data of each of the families, wherefrom is derived the formation of an orthonormal vector base (new attributes) for filtering and describing said geophysical events. The invention is applicable to geological interpretation or to interpretation of an underground reservoir, from seismic measurements or logs for example.

1. Field of the Invention
The present invention relates to a method for packing and filtering geophysical data in order to retrieve therefrom information on the nature of the subsoil.
2. Description of the Prior Art
Geophysical measurements are conventionally used to provide complementary information, in relation to direct observations in wells, on the geologic properties variations of underground formations. Seismic interpretation is for example based on the analysis of attributes in order to retrieve these geologic properties from the seismic information, in a zone of interest of an underground formation (reservoir zone for example). Another example is log interpretation where the various measurements are simultaneously analyzed to deduce therefrom the petrophysical characteristics of the underground formation studied.
A standard approach for seismic or log attributes analysis involves the use of statistical pattern recognition and estimation methods in order to calibrate the pertinent attributes with the available information obtained in various wellbores.
For example, French Patent 2,768,818 and corresponding U.S. Pat. No. 6,051,651 filed by the assignee describes a method allowing detection of the potential classes in a population of events related to the physical properties of a complex medium such as the subsoil, located from data obtained by subsoil exploration or by in-situ measurements (events of geologic, geophysical nature, events linked with fluid production data, etc.), these events being located by points in a multi-variable space and defined by an attributes vector.
Another example is given by French Patent Application 01/05,675 filed by the assignee, which describes a method for facilitating monitoring over time of the evolution of an underground zone by compared analysis of a certain number n of seismic records successively obtained after n successive 3D seismic surveys (4D seismic method), wherein a pattern recognition technique applied to the whole of the seismic events of several surveys, considered and analyzed simultaneously, is used.
EP Patent A-671,017 and corresponding U.S. Pat. No. 5,638,269 filed by the assignee can also be mentioned, which describe a method allowing establishing of a connection between geologic data obtained by coring or logging in wells and seismic data obtained by means of a seismic exploration survey, which is based on a statistical calibration technique with statistical calibration obtained by bringing together the local geologic data (measured in wells) and seismic attributes read on seismic traces obtained in the immediate vicinity of each well.
The seismic attributes that are subjected to such an interpretation are calculated in the post-stack or in the pre-stack domain.
In the post-stack domain, the conventional attributes are calculated from the amplitudes at the reservoir level, or from the impedance P estimated by means of a stratigraphic inversion.
In the pre-stack domain, the number of pertinent attributes can increase considerably. In this case, the attributes are either the amplitudes at different offsets or incidence angles, or parameters obtained from a pre-stack joint stratigraphic inversion (impedances P and S, density, product of the Lame parameters with the density, etc.).
Concerning the interpretation of logs, the situation is similar to a very large number of measurements being available, hence a large number of attributes allowing characterization of the logs, these attributes being often organized in families of equal physical significance (for example family of the resistivity attributes, family of the attributes associated with the natural or induced radioactivity of the formations, etc.).
A method referred to as principal components analysis (PCA) method, well known in the art, can be applied to these attributes in order to analyze the relations existing between them and to reduce the number of significant attributes. The principal components extracted from the PCA are the new attributes: they define an orthogonal or non-orthogonal (rotation) vector base and are linear combinations of initial variables. They can be used in statistical pattern recognition algorithms and correspond to a multivariate filtering of the initial seismic or log information. Various applications of the PCA method are for example described in:    Dumay, J., Fournier, F., 1988, "Multivariate Statistical Analyses Applied to Seismic Facies Recognition", Geophysics, 53, 1151-1159,    Hagen, D., C., 1982, "The Application of Principal Component Analysis to Seismic Data Sets", Geoexpl., 20, 93-111, or    French Patent A-2,772,138 and corresponding U.S. Pat. No. 6,345,108 filed by the assignee.
However, PCA does not take into the attributes group and, consequently, it is often difficult to give a physical interpretation of the principal components, or to relate them clearly to the initial attributes, especially if they are organized in groups of equal physical significance.
The analysis referred to as canonical analysis is also a statistical method well known in the art, which allows establishing the relations that may exist between two sets of variables in order to know if these two sets describe the same properties. This method is for example described by:    Hotelling, H., 1936, "Relations Between Two Sets of Variables", Biometrika, 28, 321-377.
An example of application of canonical analysis is for example described by:    Fournier, F., and Derain, J. F., 1995, "A Statistical Methodology for Deriving Reservoir Properties from Seismic Data", Geophysics, 60, 1437-1450.
This method is limited to the study of two sets of variables and does therefore generally not apply to the multicube seismic information or to the multidomain log information. Furthermore, it defines synthetic variables in each subspace associated with the two sets, and not a single vector base allowing packing the whole of the initial variables and to describe the relation between the two sets.
There are several methods allowing generalization of the analysis referred to as canonical analysis by extending the analysis to more than two sets of variables. Different aspects of the prior art in the field considered are described, for example, in the following publications:    Horst, P., 1961, "Relations Among M Sets of Measures", Psychometrika, 26, No. 2, 129-149,    Caroll, J. D., 1968, "A Generalization of Canonical Correlation Analysis to Three or More Sets of Variables", Proc. 76th Conv. Amer. Psych. Ass.,    Kettenring, J. R., 1971, "Canonical Analysis of Several Sets of Variables", 58, 3, 433-450, or    Saporta, G., 1990, "Probabilites, Analyse des Donnees et Statistiques", Technip, Paris.
These methods however have certain limits. In particular, they do not allow describing each set separately, the synthetic variables cannot be readily related to the various sets and, consequently, be physically interpreted. Furthermore, only the global inertia part they represent can be known, and not the inertia part they represent for each set. It is therefore not possible to filter each set independently.
Another known analysis technique referred to as generalized principal components analysis (GPCA) allows comparison of different sets of variables while describing each one. It is implemented on data of economic nature for example in:    Casin, Ph., 2001, "A Generalization of Principal Component Analysis to K Sets of Variables", Computational Statistics & Data Analysis, 35, 417-428.

Method for packing and filtering geophysical events to retrieve therefrom data on the type of subsoil


An isolation chamber and disinfection system. The novel system includes an enclosure for providing a sturdy, lightless and soundproofed environment and for containing a solution for supporting a user in a state of floatation; a disinfection system for disinfecting the solution, the disinfection system including an ultraviolet disinfection system and an ozone injection system; and a pump for removing the solution from the enclosure and circulating it through the disinfection system at a predetermined flow rate. The enclosure is formed from interlocking, foam-insulated panels and lined with a potable liner that holds the solution. In relationship to the isolation chamber, in accordance with the present teachings, an improved method of disinfection is taught using UV radiation with ozone injection. Ozone has a short half-life and leaves substantially no toxic byproducts. Verification of disinfection can be accomplished by monitoring the dissolved ozone, which can be determined accurately utilizing a dissolved ozone sensor system.

Field of the Invention
The present invention relates to isolation chambers. More specifically, the present invention relates to systems and methods for disinfecting floatation tanks.
Description of the Related Art
An isolation chamber, or floatation tank, is a dark, soundproofed enclosure containing a layer of salt water upon which a user floats, typically for relaxation or meditation purposes. The chamber is designed to isolate the user from external stimuli such as light, sound, smell, and, in particular, gravity. Epsom salt is usually added to the water in the isolation tank to increase the density of the solution such that the user floats comfortably near the surface with the face above the water. The buoyancy created by this salt solution effectively eliminates the body's specific gravity, bringing the floater close to an experience of total weightlessness.
Isolation chambers for commercial usage must maintain proper sanitation practices to ensure a safe and healthy environment for the users. As with swimming pools and other recreational water facilities, the water in an isolation chamber is subject to undesirable contaminants introduced by the users and environment. These contaminants may include pathogenic micro-organisms such as viruses, bacteria, protozoa, and fungi. Commercial isolation chambers should therefore include a properly designed disinfection system for killing these micro-organisms to prevent the transmission of infectious diseases.
Conventional isolation chambers typically rely on chemical disinfectants such as chlorine for water sanitation. Chlorine disinfection is inexpensive, quick and effective at killing most pathogenic microbes, and leaves a residual in the water that continues disinfecting for a long period of time. Residual chlorine is desirable, and often required, for public swimming pools and similar venues where a user may be exposed to contaminants introduced by other swimmers. This residual is used to validate a certain level of decontamination. This criteria is designed to be used with regular water in an open space.
In a magnesium sulfate solution one cannot accurately monitor the chemical residual, making it impossible for either the operator or public health authority to determine correct measurement of disinfecting chemicals, resulting in blind dosing.
Chemical disinfectants react with organic matter to produce toxic disinfection by-products such as trihalomethanes, haloacetic acids, chlorite, and others. In particular, chlorine reacts with ammonia or nitrogen containing compounds such as sweat, saliva, and urine to produce chloramines. Exposure to chloramines can cause respiratory problems, rashes and other skin conditions, eye irritation, digestive disorders, and other health issues. In addition, chloramines also produce an irritating odor which is particularly undesirable for isolation chambers, where the desired goal is to isolate the user as much as possible from external stimuli.
Hence, a need exists in the art for an improved system or method for sanitizing and monitoring isolation chambers that is safer than conventional chemical approaches.


A liquid jet head manufacturing method, a method for integrally manufacturing a liquid jet apparatus, a liquid jet head, a printing apparatus and a liquid jet apparatus. The manufacturing method includes: forming multiple pressure-generating members arranged spaced apart on a first substrate; forming, on a first surface of first substrate, pressure chambers corresponding to multiple pressure-generating members, and a common chamber in communication with multiple pressure chambers; forming a transition layer on pressure chambers, and forming a jet orifice plate on transition layer; and forming, on jet orifice plate and transition layer, jet orifices in communication with pressure chambers. When the number of pressure chambers needs to be increased, since pressure chambers are formed on first substrate individually, mechanical strength of first substrate will not be reduced, and during manufacturing process, breakage of first substrate can be avoided, improving yield of liquid jet head, and reducing manufacturing cost.

In terms of a liquid jet head or a liquid jet apparatus of a printer, deformation of a piezoelectric element and a vibrating plate causes a change in volume of a pressure chamber, so that ink in the pressure chamber jets from a jet orifice.
An existing liquid jet head includes a substrate, a vibrating plate and a piezoelectric element arranged on a first face of the substrate, and a jet orifice plate bonded on a second face (opposite to the first face) of the substrate. The existing liquid jet head is manufactured as follows: forming a vibrating plate and a piezoelectric element on a first face of a substrate, etching, by means of performing an etching process on a second face of the substrate, the substrate to form a plurality of pressure chambers for storing liquid corresponding to positions of piezoelectric elements and a common chamber corresponding to the position of an ink supply hole, and finally, bonding a jet orifice plate on the second face of the substrate, so that a plurality of jet orifices on the jet orifice plate are in communication with the pressure chambers respectively. When the liquid jet head is in operation, the piezoelectric element is deformed under driving of a voltage, and the deformation is transferred to the vibrating plate to cause a change in volume of the pressure chamber, so that liquid in the pressure chamber jets from the jet orifice to complete the printing.
An existing liquid jet apparatus includes a substrate, a vibrating plate and a piezoelectric element arranged on a first face of the substrate and a jet orifice plate bonded on a second face (opposite to the first face) of the substrate. The existing liquid jet apparatus is manufactured as follows: forming a vibrating plate and a piezoelectric element on a first face of a substrate, etching, by means of performing an etching process on a second face of the substrate, the substrate to form a plurality of pressure chambers for storing liquid corresponding to positions of piezoelectric elements and a common chamber corresponding to the position of an ink supply hole, and finally, bonding a jet orifice plate on the second face of the substrate, so that a plurality of jet orifices on the jet orifice plate are in communication with the pressure chambers respectively. When the liquid jet apparatus is in operation, the piezoelectric element is deformed under driving of a voltage, and the deformation is transferred to the vibrating plate to cause a change in volume of the pressure chamber, so that liquid in the pressure chamber jets from the jet orifice to complete the print.
However, in order to improve a print resolution of the printer, the number of pressure chambers of the liquid jet head or the liquid jet apparatus needs to be increased. Since the existing pressure chamber of the liquid jet head or the liquid jet apparatus is formed by etching on the substrate, increasing the number of the pressure chambers means to reduce the sidewall thickness of adjacent pressure chambers, which inevitably causes that the mechanical strength of the silicon wafer as the substrate is decreased, and the substrate is broken readily during the manufacturing process, thereby resulting in that the yield of the liquid jet head or the liquid jet apparatus is reduced and the manufacturing cost is relatively high. Besides, since the jet orifice plate is boned on the substrate of the pressure chamber by an adhesive, if the adhesive flows into the pressure chamber, print quality also can be affected.

Liquid jet head, method for integrally manufacturing a liquid jet apparatus, and device


An X-ray tube includes a cathode including an emitter emitting an electron beam, an anode at which a target material is disposed, the target material emitting an X-ray by colliding with the electron beam, and an insulating spacer isolating the anode, wherein the cathode or the anode is disposed between the emitter and the insulating spacer.

1. Field
An aspect of the present disclosure relates to a structure of an X-ray tube.
2. Description of the Related Art
FIG. 1 illustrates a general structure of an X-ray tube requiring a high acceleration voltage may be configured to include a cathode 10 for emitting an electron beam, an emitter 11, a gate 20, a focusing electrode 30, and an anode 40. The electrodes may be electrically isolated from each other by an insulating spacer 50. The insulating spacer 50 may have a tubular shape. When the emitter 11 is a thermoelectron source, the gate 20, the focusing electrode 30, and the like may be omitted. When the emitter 11 is a field emission electron, the focusing electrode 30 may be integrated with the gate 20 to have the same potential. Electrons (e−) emitted in the form of an electron beam from the emitter 11 are accelerated by a voltage difference between the anode 40 and the cathode 10 and then attracted toward the anode 40. Although not shown in this figure, when the electrons collide with a target material (not shown) disposed at the anode 40, an X-ray is emitted. The anode 40 may be an inclined anode or a transmissive anode. Since the insulating spacer 50 is positioned around the path along which the electrons accelerated with a high voltage are attracted toward the anode 40, electric charges are accumulated in the insulating spacer 50, and therefore, an abnormal operation may be caused. The electric charges accumulated in the insulating spacer 50 may be transferred to another electrode under a high-voltage atmosphere. In this case, the X-ray tube may be damaged due to flow of the electric charges in an arc form.
When the field emission electron source is used, the quantity of emitted electrons may be controlled using an active current control unit 60 configured by connecting a high-voltage field effect transistor, etc. in series to the cathode 10 as shown in FIG. 1. In this case, a reference voltage Vref of the active current control unit 60 may be a ground voltage (0V). Current limit conditions may be determined according to characteristics of a field emission emitter, gate voltages, and gate-source voltages applied to the field effect transistor. Here, the voltage of the cathode 10 may be increased as compared with the reference voltage Vref. The voltage of the cathode 10 may be fluctuated depending on a change in characteristics of the emitter 11 by the active current control unit 60 that controls a field emission current to be constant under the current limit conditions. If a gate voltage Vg, a focusing voltage Vf, and an anode voltage Va are maintained constant, focusing characteristics of an electron beam may be changed as the voltage of the cathode 10 is changed under the current limit conditions.

ELECTRONICS AND TELECOMMUNICATIONS RESEARCH INSTITUTE


A vacuum vessel that includes a first and a second substrate facing each other and extending across both an active area and a non-active area surrounding the active area, a sealing member arranged at peripheries of the first and the second substrates and adapted to maintain a vacuum between the two substrates, a plurality of wall type spacers arranged between the first and the second substrates while extending across the active area and a plurality of spacer supports arranged in the non-active area between the first and the second substrates, the plurality of spacer supports including a plurality of grooves adapted to receive the ends of respective ones of the plurality of wall type spacers, each spacer support having a height identical to or greater than a height of the plurality of wall type spacers.

1. Field of the Invention
The present invention relates to a vacuum vessel, and in particular, to a vacuum vessel which has built-in spacers for spacing first and second substrates apart from each other by a predetermined distance, and an electron emission display device using the vacuum vessel.
2. Description of the Related Art
Generally, electron emission devices are classified into those using hot cathodes as an electron emission source and those using cold cathodes as the electron emission source. There are several types of cold cathode electron emission devices, including a field emitter array (FEA) type, a metal-insulator-metal (MIM) type, a metal-insulator-semiconductor (MIS) type, and a surface conduction emitter (SCE) type.
Although the electron emission devices are differentiated by specific structure depending upon the types thereof, they all basically have electron emission regions formed on a substrate, and driving electrodes for controlling the on/off and amount of electron emission from the electron emission regions. The electron emission devices can be used as an electron emission structure for a light source, such as a backlight or an image display device.
With the typical structure of the electron emission display device using the electron emission device, electron emission regions and driving electrodes are formed on a first substrate, and phosphor layers are formed on a surface of a second substrate facing the first substrate together with an anode electrode that keeps the phosphor layers at a high potential state. The first and the second substrates are sealed together at their peripheries using a sealing member, and the interior thereof is exhausted to form a vacuum vessel so that the electrons can be fluently emitted and migrated therein. A strong compression force is applied to the vacuum vessel due to the pressure difference between the interior and exterior thereof. A plurality of spacers are provided within the vacuum vessel to prevent vacuum vessel from breaking due to the compressive force. The spacers are attached to any one of the first the second substrates using an adhesive layer, and placed within the active area along with the electron emission regions and the phosphor layers.
With such a vacuum vessel, when spacers are fitted to one of the first and the second substrates and the interior thereof is evacuated, the spacers and the other of the first and the second substrates spaced apart from each other without an intervening adhesive layer are held tightly in contact with each other so that an impact is applied to the spacers, and the spacers are liable to be broken due to this impact.
The electron emission display device further has a non-active area located between the active area and the sealing member that does not serve to display an image. With the distribution of the stresses applied to the first and the second substrates after the exhausting, the stress applied to the non-active area is greater than the stress applied to the active area. This is because the structure for absorbing and withstanding the pressure of the two substrates is not present in the non-active area. Accordingly, cracks are likely to occur in the vacuum vessel due to the relatively large stress in the non-active area.
As spacers are attached to the substrate using an adhesive layer, the adhesion thereof with respect to the substrate is relatively weak. Consequently, some spacers are inclined or detached from the substrate during the exhausting process so that the pressure applied to the vacuum vessel is not uniformly distributed. As a result, the inclined spacers can block the paths of the electron beams, thus deteriorating the display characteristic. Moreover, as the wall type spacers have a high sectional aspect ratio and a long length, they are prone to twisting. For this reason, in a vacuum vessel using the wall type spacers, the spacers are likely to be twisted or inclined after the exhausting. Therefore, what is needed is an improved design for a vacuum vessel and an electron emission display device having the same that is better able to withstand and absorb the pressure caused by the vacuum vessel while overcoming the above problems.

Vacuum vessel and electron emission display device using the same, provided with spacer supports in non-active area of the display


A pivot and slide door system for an automotive vehicle includes a door adapted to be disposed in a door opening of a vehicle body. The pivot and slide door system also includes a pivot assembly attached to the vehicle body and cooperating with the door to pivot the door from a closed position to a pivot position. The pivot and slide door system further includes a slide assembly cooperating with the pivot assembly to slide the door from the pivot position to an open position.

1. Field of the Invention
The present invention relates generally to door systems for vehicles and, more particularly, to a pivot and slide door system for an automotive vehicle.
2. Description of the Related Art
It is known to provide a door system for a vehicle such as an automotive vehicle to allow an occupant to enter and exit an occupant compartment through a door opening in a vehicle body of the automotive vehicle. The door system typically includes a door extending longitudinally to close the door opening. The door is pivotally mounted at a front end by a hinge to vehicle structure forming the door opening on the vehicle body. The door system also includes a latch at a rear end of the door to latch the door in a close position to a striker connected to the vehicle structure. In operation, an operator moves a handle to unlatch the door from the striker and pivot the door via the hinge to open the door opening. To close the door, the operator pivots the door via the hinge to close the door opening and engage the latch with the striker.
Although the above door system has worked well, it is desirable to open and close the door opening via the door when the automotive vehicle is parked close to an object. It is also desirable to improve the ingress/egress through the door opening to the occupant compartment via the door. It is further desirable to improve the ease of operation of the door independent of the size of the door. Therefore, there is a need in the art to provide an improved door system for an automotive vehicle.


A method of designing an improved bearing system is provided. In one embodiment, the method for designing a fluid dynamic bearing system includes determining a first stability ratio for a first journal bearing configuration. The method further includes determining a second stability ratio for a second journal bearing configuration. In one embodiment, the first configuration has two sub-journal bearings and the second configuration has three sub-journal bearings. The two stability ratios may then be compared to determine whether adding a sub-journal increases the stability ratio of the bearing system.

1. Field of Invention
The present invention relates to fluid dynamic bearing motors and, more specifically, to a multi-journal fluid dynamic bearing motor assembly.
2. Description of Related Art
FIG. 1 provides a perspective view of a disc drive assembly 150. In this arrangement, a plurality of discs 110′ are stacked vertically within the assembly 150, permitting additional data to be stored, read and written. The drive spindle 151 receives the central openings 105 of the respective discs 110. Separate suspension arms 156 and corresponding magnetic head assemblies 158 reside above each of the discs 110. The assembly 150 includes a cover 130 and an intermediate seal 132 for providing an air-tight system. The seal 132 and cover 130 are shown exploded away from the disc stack 110′ for clarity.
In operation, the discs 110 are rotated at high speeds about the spindle 151. As the discs 110 rotate, an air bearing slider on the head 158 causes each magnetic head 158 to be suspended relative to the rotating disc 110. The flying height of the magnetic head assembly 158 above the disc 110 is a function of the speed of rotation of the disc 110, the aerodynamic lift properties of the slider along the magnetic head assembly 158 and, in some arrangements, a biasing spring tension in the suspension arm 156.
A servo spindle 152 pivots about pivot axis 140. As the servo spindle 152 pivots, the magnetic head assembly 158 mounted at the tip of its suspension arm 156 swings through arc 142. This pivoting motion allows the magnetic head 158 to change track positions on the disc 110. The ability of the magnetic head 158 to move along the surface of the disc 110 allows it to read data residing in tracks along the magnetic layer of the disc. Each read/write head 158 generates or senses electromagnetic fields or magnetic encodings in the tracks of the magnetic disc as areas of magnetic flux. The presence or absence of flux reversals in the electromagnetic fields represents the data stored on the disc.
Fluid dynamic bearings tend to generate less vibration and non-repetitive run-out in the rotating parts of motors than ball bearings and other types of bearings. For this reason, fluid dynamic bearing motors are oftentimes used in precision-oriented electronic devices to achieve better performance. For example, using a fluid dynamic bearing motor in a magnetic disc drive, such as magnetic disc drive 150 described above in conjunction with FIG. 1, results in more precise alignment between the tracks of the discs and the read/write heads. More precise alignment, in turn, allows discs to be designed with greater track densities, thereby allowing smaller discs and/or increasing the storage capacity of the discs.
As persons skilled in the art are aware, an ongoing challenge in fluid dynamic journal bearings is balancing the tradeoff between motor performance and power consumption. For example, increasing the stiffness of the fluid dynamic journal bearings results in less vibration in a motor's rotating parts and, therefore, increased motor precision and performance. However, an increase in the stiffness of the bearings is usually accompanied by an increase in the power consumption of the motor. Therefore, there exists a need for a technique to increase the stability of a fluid dynamic bearing without increasing the amount of power consumed by the fluid dynamic bearing.

Method for designing a fluid dynamic bearing system


An optical glucose sensor chip includes a substrate, a pair of optical elements formed on a surface of the substrate for introducing light into the substrate and for emitting the light from the substrate, and a glucose sensing membrane formed on the surface of the substrate at a position between the optical elements. The sensing membrane includes a color reagent substrate, a first enzyme which oxidizes or reduces glucose, a second enzyme that generates a material which makes the color reagent substrate exhibit color by a reaction with a product obtained by oxidation or reduction of glucose, a nonionic cellulose derivative, and an ionic polymer into which a buffer is incorporated. At least one of the first and second enzymes is coated with the ionic polymer, and the color reagent substrate. The first and second enzymes, the buffer and the ionic polymer are supported by the nonionic cellulose derivative.

1. Field of the Invention
The present invention relates to an optical glucose sensor chip and a method of manufacturing the optical glucose sensor chip.
2. Description of the Related Art
As an optical glucose sensor chip, for example, a less invasive type blood sugar level measuring chip has been developed which indirectly measures a blood sugar level by extracting body fluid of subcutaneous tissues. This sensor chip has a structure including a glass substrate, a pair of gratings which are formed on a surface of the substrate and introduce light to or emit the light from the substrate, and a glucose sensing membrane which is positioned between the gratings and formed on the surface of the substrate. This glucose sensing membrane contains a color reagent substrate (for example, 3,3′5,5′-tetramethylbenzidine [TMBZ]), a first enzyme (for example, glucose oxidase [GOD]) that oxidizes or reduces glucose, a second enzyme (for example, peroxidase [POD]) generating a material that reacts with a product obtained by the oxidation or reduction of glucose to make the color reagent substrate exhibit color, and a film forming high-molecular compound (for example, a cellulose derivative such as carboxymethyl cellulose [CMC]).
When a sheet gel is disposed between the skin and the above sensing membrane to apply an electric field in a glucose sensor chip having such a structure, glucose in a subcutaneous tissue solution penetrates the gel from the skin and reaches the above sensing membrane. At this time, the above TMBZ as a color reagent substrate in the sensing membrane makes to exhibit color due to the reaction between glucose and GOD or POD. When light is made to be incident to the above substrate and to be refracted on the surface of the substrate and on one of the above gratings, the light propagates through an interface between the above substrate and the sensing membrane containing color-formed TMBZ, is refracted on the interface between the substrate and the other grating and is received by, for example, a light receiving element. The intensity of the received laser light is less than the intensity (initial intensity) of the laser light received by the light receiving element when the above sensing membrane has not exhibited color. For this reason, the concentration of the above glucose can be detected from the ratio of reduction of the above light intensity.
However, when the above sensing membrane is stored and used for a long time, the activity of the first and second enzymes in the membrane rapidly deteriorates. Examples of the cause of the deterioration include a variation in the pH of the sensing membrane, a variation in the ionic strengths of the first and second enzymes and hydrolysis of the first and second enzymes. When the first and second enzymes deteriorate, the first enzyme reacts insufficiently with glucose which is the subject to be measured. The reduction in the reactivity with glucose reduces the generation of the material that aids the color-forming material obtained by the subsequent reaction with the second enzyme in exhibiting color, with the result that this causes a reduction in the reactivity with the color-forming material and a reduction in the degree of the color to be exhibited, bringing about a deterioration in the sensitivity of the glucose sensor chip.

Optical glucose sensor chip and method of manufacturing the same


A learning mode 40 permits a device 11 to register a plurality of transmitters (as part of authorizing such transmitters to act as remote control devices for use with a movable barrier operator 12) without necessitating intervening user tactile inputs (as through a user interface keypad 33) and without unnecessarily creating a security risk with respect to the duration of the learning mode. In one embodiment, duration of the learning mode is extended in response to detecting indicia that the learning mode continues to be used.

This invention relates generally to device learning modes.
Many devices, though programmed or designed to effectuate one or more specific tasks, must nevertheless be provided with supplementary information from time to time in order to properly execute such task(s). For example a movable barrier operator (such as, for example, a garage door or gate opener) can be provided by the manufacturer with functionality that permits use of the movable barrier operator with a movable barrier, but this initially provided information will often not be sufficient to ensure complete or compatible usage. As one example, a property may use a movable barrier operator to control access to a certain area. Access may be allowed to only, say, 250 individuals. If each individual is provided with a remote control transmitter to interface with the movable barrier operator, identifying information for such transmitters will ordinarily not be known at the outset to the operator. Consequently, the operator must learn this identifying information in order to operate at desired.
Various ways are known to support such a learning process. Very typically, and still using movable barrier operators as an example, a user will press one or more buttons or keys on the device to initiate the learning mode. Additional entries may be necessary to effectuate different purposes (for example, it may be necessary for the user to enter a specific memory location or alias/identifier for each given transmitter to be registered). The user then activates a transmitter and the signals from the transmitter are received by the device while still in the learn mode. The device will then store the identifying information contained in the transmission and conclude the learn mode. Such an approach works acceptably in some circumstances but is considerably deficient in other settings. For example, if the user must register a large number of transmitters (many such systems will support hundreds of such transmitters or even more), then the above process can be tedious, time consuming, and prone to error.
At least one prior art approach attempts to ameliorate this problem by remaining in a learn mode for a predetermined period of time (such as 30 seconds) and allowing a number (such as 4) of transmitters to be registered during that period of time. While possibly acceptable for a small number of transmitters, again, such an approach presents little improvement for a user looking to program a considerably larger population. On the other hand, extending such the period of time for the learn mode in this approach is typically considered ill-advised, as any compatible transmission as received during this period of time will be accepted into the system. Therefore, extending the period of the learning window in this fashion is viewed as a significant security risk and hence, is eschewed as a viable approach.